Introduction

What is AP Calculus BC?

This is a term used by the College Board, which has an approximate college-level
analog; it refers to the material usually covered towards the end of a college Calculus I
course (roughly the B aspect) and most of the material covered in Calculus II (roughly
the C aspect). A BC course does not encompass a great deal of new material. It
extends the subjects of differentiation and integration to new functions; such as
parametric, polar, and vector functions; and introduces the concept of infinite series and
their applications to polynomial approximations. Most the material on the BC exam is
actually derived from the AB aspect of calculus. The format of the AP Calculus BC
exam is identical to that of the AP Calculus AB exam, save for the fact that the former
exam tests more material. There are three parts: a multiple choice section in which
calculator use is not permitted, a multiple choice section in which the calculator is
permitted and a free response section in which one may use the calculator for the first
half but not the second.
About this book
I have written this book for two reasons. Firstly, while books on the market
dedicated either specifically to AP Calculus AB or AB/BC are legion, one could
undoubtedly count the number of books concentrating upon AP Calculus BC on one
hand. From one point of view, this is understandable; there are simply more students
who take the AB test than the BC test. However, the dearth of concentration upon BC
material makes it quite difficult, in my opinion, to obtain information that is equal in
profundity to the AB material that is available. Secondly, while AP Calculus AB
introduces many interesting applications of differentiation and integration to the physical,
life, and social sciences, AP Calculus BC is almost completely devoid of interesting
applications. Even spinning a curve around an axis and evaluating the volume of the
solid generated is more interesting than blindly approximating values of functions using
Taylor polynomials! While I do include discussions on the pure-mathematical aspects of
the BC material, I also introduce interesting, perhaps somewhat biased, applications of
the material that is not tested on the AP exam, as the title of this book shows. While
many of the applications that I discuss are not tested nationally, I believe that they foster
a higher degree of appreciation for the BC material. I apologize to the student
economists and social scientists using this book; those areas are not my specialty.
Nevertheless, I feel that my combination of pure mathematics and the physical and life
sciences in this book will nurture a greater interest in the material and, perhaps, even a
sense of freedom from the burden of the standardized test.
I would also like to comment on the last chapter of this book, entitled Vector
Calculus and Curves in Space. The AP Calculus BC exam only tests a very small
portion of the material that I include, mainly differentiating and integrating vector-valued
functions. My discussion goes well beyond this for a good reason; vector calculus is an
excellent bridge into Calculus III, which is multivariable calculus (i.e. partial derivatives,
1

multiple integrals, etc.). I hope that concluding my book in this manner will encourage
the readers research on multivariable calculus, even if he or she does not plan to take the
course in college.
Before reading the book
I have generally listed the material with which the student should be familiar
before he or she begins reading the book:
I.) Limits and Continuity: This includes understanding what a limit is and how to
evaluate one algebraically and graphically. One should also understand how tell whether
or not a function is continuous.
II.) Derivatives: One should know the limit definition of the derivative (difference
quotient) and the various rules for finding derivatives (i.e. the power rule, the product
rule, the quotient rule, and the chain rule). One should also know how to take the
derivative of all trigonometric functions, exponential functions, and logarithmic
functions. One should also understand implicit differentiation.
III.) Applications of Derivatives: It is expected that one knows how to find the slope of
the tangent line and the normal line to the curve. One must be able to compare the graph
of a function to the graph of its first and second derivatives and vice versa. One must be
able to find local minima and maxima as well as points of inflection by using derivatives.
One should understand the Mean Value Theorem and Rolles Theorem. One should be
able to do optimization, motion, and related rates problems.
IV.) Integrals: One must know the various rules for antidifferentiation (i.e. the power
rule, u-substitution, the natural logarithm, inverse trigonometric functions). One should
know how to approximate the area under a curve by using rectangles and the Trapezoidal
Rule and get the exact area under a curve and between curves by using the 2nd
Fundamental Theorem of Calculus. One should also be familiar with the 1st Fundamental
Theorem of Calculus and accumulation functions.
VI.) Applications of Integrals: It is expected that one knows how to apply the integral to
problems in motion. One must know how to find the volume of a solid of revolution by
using the Disc Method, the Washer Method, the Shell Method, and various crosssections. One should know how to evaluate the average value of a function.
VII.) Differential Equations: One must be able to solve differential equations via
separation of variables. One should know how differential equations apply to
exponential change and Newtons Law of Cooling. One should also be able to
understand slope fields and linearization.

A Note on Technology:
Nowadays in advanced mathematics courses, the graphing calculator is used
extensively. While this machine is incredibly useful, it is important not to let it overcome
the power of the mind. Before using the calculator to solve a problem graphically or
analytically, the student should understand the theory behind the steps that the calculator
took to solve the problem. It is also important to remember that half of the AP Calculus
BC test prohibits the use of a calculator. Therefore, wherever possible, the student
should do calculus in the 19th-century fashion (even without a slide rule!) by hand. As
of the completion of the manuscript of this book, the TI-83 graphing calculators have
been discontinued. Thus, when discussing solutions via calculator, I will refer to its
contemporary counterpart, the TI-84. I would also like to mention that throughout the
book, notably in the final chapter (Vector Calculus), I have generated some of the graphs
using the computer program Mathematica. This program is one of the most powerful
computational tool in the world, and I encourage readers to learn more about its nearly
limitless capabilities.

Table of Contents
Prelude: At the Level of the Infinitesimal.5-7
Chapter 1: LHpitals Rule and Advanced Techniques
of Integration.8-24
Chapter 2: Differential Equations. 25-39
Chapter 3: Infinite Sequences and Series..40-55
Chapter 4: Power Series and Polynomial Approximations...56-74
Chapter 5: New Coordinate Systems: Parametric and Polar.75-102
Chapter 6: Vectors and Vector Calculus......103-131
Practice Test with Answers..132-188
Appendix A (Essential Pre-Calculus Information)..189-190
Appendix B (Brief Table of Integrals)..191
References.192
Index193-195

Prelude: At the Level of the Infinitesimal

In modern vernacular, the word infinity is thoroughly abused, most likely the
result of a misconception of its mathematical significance. When one wishes to express
an extreme, infinitely is often inserted for dramatic effect: It was infinitely
informative, This performance was infinitely better than the last one, etc. What is the
true meaning of infinity? An attempt to imagine a physical infinity will yield no greater
victory than an attempt to conceive of more than three dimensions. Mathematics,
however, has the power not only to make sense of these abstractions, but also to put them
to good use. Differential and integral calculus is based upon the concept of infinity.
When one evaluates the derivative of a function, he or she is essentially finding the slope
of a secant line on a curve, which intersects two points on that curve, when the difference
between those two points is infinitely small. In effect, by moving from a discrete
difference between points to an infinitely small difference, the secant line has become a
tangent line! The concept of the integral also exploits the power of infinity. When one
uses the integral to find the area under a curve, for instance, he or she is actually taking
the sum of an infinite number of infinitely small chunks of an area to yield the entire
area underneath the curve.

Why on Earth is any of this useful? An appreciation for the practicality of

differential and integral calculus requires a little philosophical experimentation. A story
often discussed in the introduction to calculus courses is that of Zenos paradoxes. Zeno
of Elea was an ancient Greek philosopher who introduced a philosophical problem that
was largely left unsolved before the advent of calculus. According to Zenos paradoxes,
motion is illusory and impossible. While there are some variations on the specifics of his
argument, it addresses three principle dilemmas. Firstly, if a particularly fast runner
condescends to give a slower runner a head start, it will always be the case that the slower
runner will win the race. Why? When the fast runner arrives at the spot at which the
slower runner began, the latter has already moved a certain, albeit short, distance. It will
then take the fast runner a certain period of time to reach that short distance, while the
slower runner has advanced further, ad nauseam. Even more striking is Zenos proposal
that motion is futile. If a body is to move from point A to point B, which are separated
by a length l, that body must first move a distance of 1/2l. Before that, the body must
move a distance of 1/4l. Before that, 1/16l. Before that 1/32l, ad nauseam again. In fact,
not only is it the case that motion is futile, it is also impossible. According to Zeno, at
each discrete instance of time, a body is momentarily at rest. If this is the case, then there
is no point during an objects journey that it could be considered in motion. Of course,
motion is possible and a faster runner can overtake a slower runner despite a head start
5

for the latter. Therefore, there must be a solution to Zenos paradoxes. This is where
calculus makes its entrance. While the mathematical philosophy of moving to an
infinitesimal level was manifest in antiquity, calculus did not emerge as a serious
discipline and practical tool until the early 18th century, when scientists were rigorously
attempting to find solutions to difficult problems in physics. Two prominent figures
during this period, Isaac Newton and Gottfried Leibniz, are usually credited with the
invention of differential and integral calculus. While there are some philosophical
naysayers who believe otherwise, calculus very concisely solves Zenos paradoxes.
Through calculus, it can be shown that even an infinite number of distances can have a
finite sum. Furthermore, it is not the case that an infinite span of time would be required
to travel an infinite number of infinitely small distances. Infinitesimal distance (dx) and
infinitesimal time (dt) have finite significance in differentiation and integration. For
dx
instance, the expression
= v , which means that the derivative of the position function
dt
is equal to the velocity function assumes a very helpful form when integrated:
dx = vdt , which means that an infinite sum of infinitely small distances yields the
infinite sum of products of instantaneous velocity and infinitely small time spans. This
relates infinitely small changes in position with infinitely small changes in time, so,
despite the notion of the infinitesimal, motion does exist!
It is often necessary in the sciences and business to move to the infinitesimal, or
differential, level. Why? Without doing so, one is restricted to differences between two
points and ignores what is in between! In business, for example, it is necessary to
optimize the volume of production to maximize profit. This is achieved by finding the
number of products at which the difference between revenue and cost of manufacture
(profit) is a maximum. Two isolated points for revenue and cost is not sufficient; an
infinite number of such points is needed! One must evaluate the derivative of the profit
function and find the point at which it undergoes a transition from positive to negative
values, a method covered in AP Calculus AB. In the pure and applied sciences, moving
to the differential level is indispensable in modeling phenomena. Similar to the business
example, calculus allows one to consider all of the points that model an event, not just
two isolated points. Furthermore, analysis of a differential aspect of a discrete body
easily allows for the extension to the whole body; that is, through the methods of
calculus, the mathematical equations that describe an infinitely small sliver of an object
also describe that object as a whole. This somewhat bottom-up approach is very often
more practical that analyzing the object as a whole. In physics, for instance, objects have
a property known as the center of mass, the point at which the mass of the object can be
considered concentrated. In effect, one can consider infinitely small chunks of an
object, each having a mass mi. When the average of these masses is found and weighted
(statistically) by each masss position with respect to a reference point, the center of mass
i=n

m r
Mathematically, center of mass = lim
m
i =1

i i

rdm

, where ri is the position

M
of the ith infinitesimal mass with respect to a reference point and M is the total mass of
the object.

is found.

In this book, the theme of moving to the differential level will continue to be
emphasized as a useful method for understanding both the pure-mathematical aspect of
AP Calculus BC topics and their applications to the sciences.

Chapter 1: LHpitals Rule, Advanced

Techniques of Integration, and Improper
Integrals
Since AP Calculus AB usually begins with an introduction to limits and their
algebra, this book for AP Calculus BC will begin with a further discussion on limits.
Oftentimes, one cannot evaluate a limit via the methods prescribed in AP Calculus AB.
For instance, these methods do not suffice when they result in an indeterminate form, a
mathematical construction that has no immediate arithmetic meaning. There are seven
0
instances of these forms: , ,0 ,1 ,0 0 , 0 , and . Each will be discussed
0
separately.
LHpitals Rule
LHpitals Rule was developed by the Frenchman Guillaume de LHpital
(pronounced loh-pee-tahl) during the 17th century. He published this rule in his
L'Analyse des Infiniment Petits pour l'Intelligence des Lignes Courbes (1696), which is
regarded as the worlds first textbook on differential calculus! LHpitals Rule can be
expressed as follows:
If lim f ( x) = A and lim g ( x) = B , and A and B are either both equal to infinity or
x a

x a

both equal to zero, the following holds:

f ( x)
f ' ( x)
f ' ' ( x)
f n ( x)
= lim
= lim
= lim n
lim
x a g ( x)
x a g ' ( x)
x a g ' ' ( x)
x a g ( x)
The above relation states that if the ratio of two functions has a limit as x
approaches a of indeterminate form, the derivatives of both the numerator and the
denominator can be found until the limit has a determinate form. While it will not be
presented here, a proof of LHpitals Rule can be found from Cauchys Mean Value
Theorem.
There are certain limitations of LHpitals Rule that one must take into account.
0

Firstly, the rule can only be used for limits in the indeterminate forms of and ; one
0

must not attempt to use LHpitals Rule to evaluate limits that are not in indeterminate
form. Thus, it is imperative that limits be first evaluated by conventional methods and
1 1
then, only if the need arises, be subject to the rule. Note also that the forms , , and
0

are not indeterminate; they have meaningful arithmetic significance. The first
0
fraction yields infinity, the second fraction yields zero, and the third fraction yields
infinity.

0
results,
0
the derivative of the numerator and denominator may be taken until a determinate form is
achieved.
The Indeterminate Form 0/0: If one attempts to take the limit of a quotient and

tan x
.
x 0
x
tan x tan(0) 0
= Use LHpitals Rule
Solution: lim
=
x 0
x
(0)
0

Ex.) Evaluate lim

sec 2 x
sec 2 (0)
tan x
= lim
= lim
= 1.
x 0
x 0
x 0
1
1
x

lim

results upon attempting to determine the

limit of a quotient, a similar procedure can be used.

The Indeterminate Form /: If the form

ln x
.
x x 3

Ex.) Evaluate lim

ln x ln()
=
= Use LHpitals Rule.
x3
( ) 3
1
1
1
ln x
lim 3 = lim x 2 = lim 3 =
= 0.
x x
x 3 x
x 3 x
3() 3

Solution: lim
x

In the introductory discussion of LHpitals Rule, it was stated that the rule can

0
only be used to evaluate limits of the forms and . What of the other five
0

indeterminate forms? In order to apply LHpitals Rule to these, a little algebraic

manipulation is required.
The Indeterminate Form 0: This form can be algebraically transformed into either of
the forms already discussed by dividing one by the reciprocal of the other:
1

0 0
0 = 0 =
and 0 =
= .
1 0

Note that evaluating the product of two numbers and dividing one number by the
reciprocal of the other number are different operations that mean the same thing
3 2
arithmetically! For instance, 2 3 = = = 6 . Thus, while the above technique does
1 1
2 3

change the representation of the indeterminate form 0, it does not change the
arithmetic result.
Ex.) Evaluate lim x 2 e x .
x

Solution: lim x 2 e x = () 2 e ( ) = 0
x

x2
x 2 ( ) 2
Take the reciprocal: lim x e = lim
= lim x = ( ) = Use LHpitals
x
x 1
x e

e
x
e
2
2 x 2()
x
Rule: lim x = lim x = ( ) = Use LHpitals Rule again:
x e
x e

e
2x
2
2
2
lim x = lim x = ( ) = = 0 .
x e
x e

e
2

The Indeterminate Forms 1, 00, and 0: Since these indeterminate forms involve
exponents, a logarithm must somehow be applied for a reversal. If one takes the
natural logarithm of these forms, one can then deal with the problem in the context of the
techniques already discussed. Recall the property of logarithms that ln a b = b ln a . Once
the value of the limit is determined after taking the natural logarithm and applying the
previous techniques, it is important to remember to exponentiate (take the e of) of the
answer to know the original functions behavior.

Ex.) Evaluate lim+ x

x 1

1
x 1

(Note that the + subscript denotes a right-hand limit, which, along with lefthand limits, is discussed in AP Calculus AB. This is a necessary specification, since one
cannot approach a value of 1 from the left because the function does not exist to the left!)
lim x

1
x 1

x 1+

= (1)

1
(1) 1

= 1 Take the natural logarithm of the function:

1
1
ln(1) = 0 Take the reciprocal:
lim+ x x 1 = lim+
ln x =
x 1 x 1
x 1
(1) 1
1
1

(1) 1
x 1
1

= Use LHpitals Rule:

lim
= lim+
ln x = lim+
x 1+ x 1
x 1
x 1
1
1

ln x
ln(1)
1

x 1 ( x 1) (1 1) 0

=
=
= Use LHpitals Rule again:
lim
x 1+
ln(1) 0
1
ln x

ln x
1
x

10

1
1
1
1 1
= lim+
=
=
=
= . Note
x 1
0
0
1
x 1 ln x ln(1)
x ln x
2 2
1
x (1)

x 2

that is not the final answer. Since the problem was manipulated through the use of
the natural logarithm, this must be undone through exponentiation. Thus, the final
answer is e = 0.
The Indeterminate Form : Transforming this form into one that is workable
requires an algebraic manipulation of terms to yield two terms such that either one or
both are fractions. Once this is accomplished, one can multiply each term by a common
denominator to trivialize the subtraction. Since this explanation is undoubtedly not as
clear as the others, it is important to analyze the following example closely..
1
1

Ex.) Evaluate lim
x 0 sin x
x

In order to make the operation of subtraction obsolete, the first fraction is

multiplied by x and the second fraction is multiplied by sin x:
1
1
1
1
lim

= Use the technique discussed:

=
x 0 sin x
x sin(0) (0)

1
x sin x (0) sin(0) 0
1
lim
=
= Use LHpitals Rule:
= lim
x 0 sin x
x x 0 x sin x
(0) sin(0)
0

x sin x
1 cos x
1 cos(0)
0
lim
= lim
= lim
= Use LHpitals
x 0 x sin x
x 0 x cos x + sin x
x 0 (0) cos(0) + sin(0)
0
Rule again:
1 cos x
sin x
sin(0)
0
lim
= lim
=
= = 0.
x 0 x cos x + sin x
x 0 x sin x + cos x + cos x
(0) sin(0) + 2 cos(0) 2
lim+

( x 1)
= lim+
x 1
ln x

Advanced Techniques of Integration

AP Calculus AB introduced the concept of the integral and several general
techniques of integration, such as u-substitution. While these skills are quite valuable,
they are often of no use when confronted with more difficult problems, such as those that
arise in the real world. In fact, the techniques of integration introduced in this book do
not even scratch the surface of the myriad of techniques that exist. Indeed, integration is
much more difficult, and requires more rigorous mathematics, than differentiation partly
due to the fact that one is working backwards. In this section two new techniques of
integration will be covered: integration by parts and integration by partial fractions.
Integration by Parts: Integration by parts was developed by the 18th-century
English mathematician Brook Taylor, for whom Taylor Series are named, which are
studied in Chapter Four of this book. The technique can be described as the integral
analog of the product rule used in differentiation. In fact, the formula for integration by
parts is derived from the product rule for differentiation:

11

d
dv
du
uv = u + v
d (uv) = udv + vdu
dx
dx
dx

d (uv) = udv + vdu uv = udv + vdu

The last equation can be rearranged into the form: udv = uv vdu . This equation
allows one to decompose an integral into two manageable parts. The most renown
example of the necessity for integration by parts lies in the integration of the natural
logarithm.
Ex.) Evaluate

ln xdx .

The form of this integral is deceptively simple. While it may seem that the
integration of this function is relatively straightforward, the techniques of AP Calculus
AB cannot be used here. Instead one must use the technique of integration by parts. One
dx
du
(because
, in
chooses u to represent ln x and dv to represent dx. In this way, du =
x
dx
this case, is the derivative of ln x) and v = x. Now the formula for integration by parts
dx
= (ln x)( x) dx = x ln x 1 + C .
can be applied: ln xdx = (ln x)( x) x
x
In this last problem, it was relatively simple to determine what to choose as u and
what to choose as v. This determination, however, is not always so clear. One should not
blindly choose these values, for making an incorrect choice could cost a great deal of
time and energy. While it may be intuitive to the reader which functions to choose, it is
easier to follow a heuristic (rule of thumb), conveyed by the acronym LIPET
(Logarithmic functions, Inverse trigonometric functions, Polynomials, Exponential
functions, and Trigonometric functions). Basically, if one encounters an integral that
must be solved via integration by parts, he or she should first look for a logarithmic
function to represent u. If a logarithmic function is not present, one should choose an
inverse trigonometric function to represent u, and so on, in accordance with the LIPET
algorithm. Also note that integration by parts may need to be performed several times

before a manageable integral is achieved. That is, the

vdu term might require the

technique to be performed on it. When integration by parts must be performed on

multiple occasions, it is imperative not to switch back and forth between choices of
functions for u and dv. For instance, if x 3 is chosen as u in the first round of integration
by parts, 3x 2 should be chosen as u in the second round. Otherwise, one would wind up
undoing the first step!
Ex.) Evaluate

2 x 2 e 2 x dx .

The integrand contains a polynomial and an exponential function. In accordance

with LIPET, 2x 2 is chosen as u and e 2 x dx is chosen as dv. Therefore, du = 4 xdx and
1
v = e 2 x . While it should be obvious in the context of the material of AP Calculus
2
AB, v was obtained through integration by u-substitution (Remember that

12

 e du = e
u

+ C. In this case, u = 2 x and du = 2dx. Therefore, while the integrand

must be multiplied by 2 to obtain the form e u du , the integral must be multiplied by 1/2 in order to compensate for this algebraic manipulation. Again, this is AP Calculus
AB material).
1
1
1
1 2 x
2 x 2 e 2 x dx = (2 x 2 )( e 2 x ) e 2 x 4 xdx = (2 x 2 )( e 2 x ) +
e 4 xdx.
2
2
2
2
The integral is still not in a manageable form. Integration by parts must be
performed upon this last integral. In this case, u = 4 x , dv = e 2 x dx , du = 4dx , and
1
v = e 2 x .
2

2 x

4 xdx =

1
1
1
1
(4 x)( e 2 x ) e 2 x 4dx = (4 x)( e 2 x ) e 2 x 2dx = (4 x)( e 2 x ) e 2 x + C .
2
2
2
2
The last phrase is not the final answer. Recall that this is the solution to the

integral in the first round of integration by parts. That is, it is the solution to e 2 x 4 xdx .
1
The solution to this integral must be multiplied by 1/2 and added to (2 x 2 )( e 2 x ) . The
2
final answer is determined in the following way:

2 x 2 e 2 x dx =

1
1
1
1
1

= (2 x 2 )( e 2 x ) +
e 2 x 4 xdx = (2 x 2 )( e 2 x ) + (4 x) e 2 x e 2 x + C =
2
2
2
2
2

1
1
1
x 2 e 2 x + e 2 x e 2 x C = x 2 e 2 x + e 2 x C . This is the final answer.
2
2
2
Note that, since C is merely an unknown constant, 1/2C is merely another C.
Some integration by parts problems require addition of integral expressions, as the
following example shows:
Ex.) Evaluate

e cos xdx.
x

Using LIPET, e x is chosen as u and cos xdx is chosen as dv. Thus, du = e x dx

and v = sin x.

e cos xdx = (e )(sin x) (sin x)(e )dx

x

Another round of integration by parts must be performed on the new integral,

where u = e x , dv = sin xdx , du = e x dx , and v = cos x

(sin x)(e )dx = (e )( cos x) ( cos x)(e )dx . Note that this last
x

mathematical phrase is the solution to the integral yielded by the first round of integration
13

by parts. That is, it is the solution to

(sin x)(e )dx. Thus, this solution appears in the

x

solution to the overall (original) integral:

e x cos xdx = (e x )(sin x) (e x )( cos x) ( cos x)(e x )dx

e cos xdx = (e )(sin x) + (e )(cos x) e cos xdx .

x

Notice that the integral on the far right is the same as the integral that one wishes
to solve. This is where addition comes in:

e cos xdx = (e )(sin x) + (e )(cos x) e cos xdx

+ e cos xdx
+ e cos xdx

e (sin x + cos x)
2 e cos xdx = (e )(sin x) + (e )(cos x) e cos xdx =
+C.

2
x

Integration by Parts Using the Tabular Method: While the previous problems were
undoubtedly mathematically tedious, imagine having to perform integration by parts
three, four, or more times for a single problem! There is, in fact, a method for integrating
by parts that is far more elegant (and easy!). This is known as the tabular method. The
following algorithm explains how to employ this useful technique:
1.) Make a table with three columns, labeled u, dv, and 1.
2.) Choose a value for u such that it will eventually yield zero when differentiated
enough (e.g. polynomials), put it in the first row of the u column and differentiate down
the column until zero is reached.
3.) Choose a value for dv and put it in the first row of the dv column and integrate
down the column. Constants of integration are not necessary here.
4.) Put a +1 in the first row of the 1 column and alternate between positive and
negative signs down the length of the column.
5.) Draw a diagonal line from each row in the u column such that it passes
through the appropriate row in the dv and 1 columns. Once the zero in the u column is
reached, cease drawing the diagonal.
6.) Multiply all three terms connected by a given diagonal and add all of these
multiplied terms to yield the final answer.
The algorithm sounds confusing, but it usually takes less than a minute to
compute the integral. The tabular method is best conveyed through an example:
Ex.) Evaluate

x 5 cos xdx. .

The polynomial ( x 5 ) will be chosen as u and cos xdx will be chosen as dv. With
these values, one may generate the table to solve the integral:

14

u
x5
5x 4
20x 3
60x 2
120 x
120
0

dv
cos x
sin x
cos x
sin x
cos x
sin x
cos x

1
+1
-1
+1
-1
+1
-1
+1
-1

Notice that the first diagonal connects the terms x 5 , sin x , and -1; the second diagonal
connects the terms 5 x 4 , cos x, and 1 ; the third diagonal connects the terms
20 x 3 , sin x, and +1; and so on. The terms of each diagonals are multiplied and these
multiplied terms are then added:

x cos xdx. =
5

x 5 sin x + 5 x 4 cos x 20 x 3 sin x 60 x 2 cos x + 120 x sin x + 120 cos x + C . The length of
this answer indicates that solving the problem without the use of the tabular method is
quite rigorous. Indeed, the calculation takes up a couple of pages if one uses the formula
for integration by parts alone! Note, however, that the tabular method can only be used
when the u term is not infinitely differentiable, i.e., when it is a polynomial.
Integration by Partial Fractions: Upon adding fractions with dissimilar
denominators, one looks for common denominators to yield one equivalent fraction. For
x
2x
(2 x)(4 x + 2) + (3x + 1)( x)
. This is elementary algebra.
instance,
+
=
3x + 1 4 x + 2
(3 x + 1)(4 x + 2)
Working backwards is slightly more difficult, and can often be much more difficult. How
can one decompose a fraction into a sum of fractions? To do this, one would use the
method of partial fraction decomposition. In this method, a fraction is decomposed
into a sum of fractions with denominators that are the factors of the denominator of the
original fraction and with numerators that contain unknown constants. For instance, the
x 2 + 3x + 2
, has a denominator with factors of x and x 2 + 2 x + 1 (i.e.
fraction, 3
2
x + 2x + x
3
x + 2 x + x = x( x 2 + 2 x + 1) ). Thus the decomposition is as
x 2 + 3x + 2
A
Bx + C
, where A, B, and C are unknown constants.
= + 2
3
2
x + 2x + x x x + 2x + 1
Notice how the degree of the numerator is one less than the degree of the denominator.
In the denominator of the first fraction, x is linear, so the numerator must be of zero
order, with a constant of A. In the denominator of the second fraction, x 2 + 2 x + 1 is
quadratic, so the numerator must be linear, with Bx+C. Note that it is possible for factors
( x + 2)
of the denominator to be repeated, as in the fraction
, which is decomposed as,
( x + 1) 5
follows:

15

( x + 2)
A
B
C
D
E
. Here, x+1 is known as the
=
+
+
+
+
5
2
3
4
( x + 1) ( x + 1)
( x + 1)
( x + 1)
( x + 1)
( x + 1) 5
repeated factor, and repeats five times in the decomposition process, with the value of
the exponent increasing from 5 to 1. Notice that the denominators are actually linear;
this is why all of the numerators are constants. The denominators are still linear despite
their being raised to a power. This is simply a repetition of linear factors. This
discussion of partial fractions could be proven by a corollary to the Fundamental
Theorem of Algebra. However, since this is not an advanced algebra textbook, it will not
be presented here.
Where does integration play a role in all of this? Similar to the philosophy of
integration by parts, certain expressions that are not integrable by conventional methods
can be broken down into more accessible components. In the method of integration by
partial fractions, a fraction that is not integrable through the use of the methods already
learned is decomposed into smaller fractions that can be integrated through the use of
conventional methods. Remember that these fractions contain unknown constants that
must be determined. This determination is first carried out by multiplying both sides of
the equation (i.e. the whole fraction side and the partial fractions side) by the
denominator of the whole fraction. This operation cancels the denominators on both
sides of the equation. When this is completed, one can solve for a system of equations.
x 2 + 3x + 2
dx .
Ex.) Evaluate
x3 + 2x 2 + x
x 2 + 3x + 2
A
Bx + C
= + 2
3
2
x + 2x + x x x + 2x + 1
Multiply both sides of the equation by x 3 + 2 x 2 + x :
x 2 + 3x + 2
Bx + C
A
= ( x 3 + 2 x 2 + x) + 2
( x 3 + 2 x 2 + x) 3
Remember that
2
x + 2x + x
x x + 2x + 1
x 3 + 2 x 2 + x = x( x 2 + 2 x + 1) .
x 2 + 3x + 2 = A( x 2 + 2 x + 1) + ( Bx + C ) x

x 2 + 3x + 2 = Ax 2 + 2 Ax + A + Bx 2 + Cx
x 2 + 3x + 2 = ( A + B) x 2 + (2 A + C ) x + A
This last equation is a standpoint from which one can solve a system of equations. Since
the coefficients on both sides of the equation must be equal,
( A + B) = 1
(2 A + C ) = 3
A=2
This simple system of equations is rather easy to solve. Since the constant A is now
known, this value can be substituted into the other equations that contain A to find values
for B and C. When this is done, it is determined that B = -1 and C = -1. With the
decomposition completed, the integral can be easily solved:
x 2 + 3x + 2
2
( x 1)
dx =
dx = 2 ln x 2 ln x 2 + 2 x + 1 + C .
dx +
3
2
2
x
x + 2x + x
x + 2x + 1

16

The system of equations in the previous problem was not at all daunting.
Oftentimes, however, the system of equations is so horrendous that one must use matrix
algebra to determine the unknown constants, which is not part of the AP Calculus BC
curriculum. In the field of mathematics known as linear algebra, matrices, arrays of
numbers that have a myriad of uses, can be used to solve a complex system of equations.
While a thorough discussion of matrix algebra will not be presented here, this point in the
book allows an excellent opportunity to use a graphing calculator to solve difficult
problems. In the case of a complex system of equations, the TI-84 graphing calculator
can represent a system of equations as a matrix and can then transform this matrix into
one of reduced row echelon form (RREF). A matrix in RREF has the following
properties: each non-zero row has a greater number of leading zeros than the previous
row, the first non-zero number in a row is 1, and the initial 1 of row is the only nonzero element in the column in which the 1 appears. Observe the matrix below to
clarify these requirements:
1 0 0 0 2

0 1 0 0 4
0 0 1 0 3

0 0 0 1 9

Notice how each row has one more has one more leading zero than the previous
one, though it is not a requirement that each row differ by just one zero. Note also that
whenever a value of 1 appears for the first time in a row, there are only zeros in the
column in which the 1 appears. Through the methods of linear algebra, namely a
method known as Gaussian elimination, a matrix that represents a system of equations
can be reduced to RREF so that the elements in the final column represent the values of
the unknowns. To understand how this is done, see the following example. Again, this
technique is not tested on the AP Calculus BC exam, but is very useful nonetheless.
3 x 4 + 4 x 3 + 16 x 2 + 20 x + 9
as the sum of three integrals.
Ex.) Represent
( x + 2)( x 2 + 3) 2
The first step is to decompose the fraction into a sum of partial fractions. Note
that the term x 2 + 3 is a repeated factor.
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9
A
Bx + C
Dx + E
=
+ 2
+ 2
2
2
( x + 2) ( x + 3) ( x + 3) 2
( x + 2)( x + 3)

3x 4 + 4 x 3 + 16 x 2 + 20 x + 9
( x + 2)( x + 3)
=
( x + 2)( x 2 + 3) 2
A
Bx + C
Dx + E
=
+ 2
+ 2
( x + 2)( x 2 + 3) 2
( x + 2) ( x + 3) ( x + 3) 2
2

3x 4 + 4 x 3 + 16 x 2 + 20 x + 9 = A( x 2 + 3) 2 + ( Bx + C )( x + 2)( x 2 + 3) + ( Dx + E )( x + 2)
3x 4 + 4 x 3 + 16 x 2 + 20 x + 9 = Ax 4 + 6 Ax 2 + 9 A + Bx 4 + 3Bx 2 + 2 Bx 3 + 6 Bx + Cx 3 + 3Cx + 2Cx 2 + 6C
+ Dx 2 + 2 Dx + Ex + 2 E

17

3x 4 + 4 x 3 + 16 x 2 + 20 x + 9 = x 4 ( A + B) + x 3 (2 B + C ) + x 2 (6 A + 3B + 2C + D) + x(6 B + 3C + 2 D + E )
+ (9 A + 6C + 2 E )
A+ B = 3
2B + C = 4
6 A + 3B + 2C + D = 16
6 B + 3C + 2 D + E = 20
9 A + 6C + 2 E = 9
This is the system of equations that must be solved through matrix algebra. The
matrix will have as many rows as there are powers of x, including the power of zero for
the constant, and as many columns as there are unknowns plus one for the equality.
Since the system of equations in this example is associated with a fourth-order
polynomial, there will be five rows, and since there are five unknown constants, there
will be six columns. After the matrix is set up, each row is treated like one of the
equations in the system. For example, A + B = 3 , in matrix language, appears as
(1 1 0 0 0 3) , because the coefficient of A is one, the coefficient of B is one, there
are no Cs, Ds, or Es, and the whole equation is equal to 3. The final matrix is:
1 1 0 0 0 3

0 2 1 0 0 4
6 3 2 1 0 16

0 6 3 2 1 20
9 0 6 0 2 9

This matrix can be set up in the TI-84 calculator as follows: Go to 2ND + x-1,
which will open up the MATRIX menu. Go to EDIT and choose a matrix to edit (e.g.
[A], [B], etc.). The main screen of the calculator will ask for the number of rows and
columns. For this example, in put 56. Fill in the matrix.
This matrix must now be reduced to RREF. To do this via calculator, go to the
MATRIX menu again and go to the MATH option. Scroll down to B: rref( and select
this option. On the main screen of the calculator will appear the operation rref(. Input
the matrix (either [A], or [B], etc.) into this operator by choosing it from the MATRIX
menu. Press the ENTER key to calculate the RREF matrix. For this example, it should
1 0 0 0 0 1

0 1 0 0 0 2
be: 0 0 1 0 0 0

0 0 0 1 0 4
0 0 0 0 1 0

A =1
B=2
Therefore, C = 0
D=4
E=0
1
2x
4x
.
and the partial fractions are
+ 2
+ 2
( x + 2) ( x + 3) ( x + 3) 2

18

The problem merely asked for the setting up of the integral as a sum of integrals. Thus,
the final answer is:
3 x 4 + 4 x 3 + 16 x 2 + 20 x + 9
=
( x + 2)( x 2 + 3) 2
0.60976
1.19512 x + 1.60976
5.53659 x 3.07317
=
+
+
.
2
( x + 2)
( x + 3)
( x 2 + 3) 2

Improper Integrals
b

Recall the Second Fundamental Theorem of Calculus:

f ( x)dx = F (a) F (b) .

a

This theorem states that a definite integral is evaluated by finding the antiderivative of a
function, substituting the lower and upper limits of integration into this expression, and
subtracting the substitution of the upper limit from the substitution of the lower limit. In
the cases in which this theorem is employed, f(x) is bounded between a and b and the
function is continuous between and at these points (i.e. on a closed interval). A definite
integral with these properties is known as a proper integral. Whenever one or both of
these properties is not fulfilled, the expression is an improper integral. For instance, the
integral

+
2

x dx is an improper integral because the upper limit of integration is infinity.

0

Therefore, there are no restrictions on the interval. In another case, both limits may be
finite, but the function could exhibit a discontinuity within the interval in which
+1
1
dx is an improper integral
integration is to take place. For instance, the integral
1 x
because there exists an infinite discontinuity at x = 0. Interestingly, it may very well be
the case that an improper integral represents a finite area despite an infinite interval or an
infinite discontinuity. Such improper integrals are said to be convergent. Improper
integrals that do not represent a finite area are said to be divergent. In order to apply the
Second Fundamental Theorem of Calculus to improper integrals, it is necessary to define
the limit(s) of integration or the value of discontinuity that makes for impropriety as a
limit. This is the major theme of calculus yet again thinking on an infinitesimal level.
Indeed, it is not possible to think of an infinite geometric representation as having a finite
area. Nevertheless, by applying the definition of the limit to the problem, it is very much
possible!
1
from 1 x < finite? If so, calculate the
Ex.) Is the area under the curve y =
3x
value of this area.
+
1
dx
What the question is essentially asking is Does the integral
3x
1
converge?
To determine this, it is necessary to define the upper limit of integration as a limit so that
the Second Fundamental Theorem of Calculus can be used. Remember, infinity is not a
number, so it technically cannot be used in this theorem. Thus, the area (A) is defined as:

19


A = lim

A = lim

a +

a
1

1
dx , where a is some value that approaches infinity.
3x

1
1
1
dx = lim ln x 1a = lim [ln(a ) ln(1)] = 0 =
a + 1 3 x
3 a +
3 x+
Therefore, no finite area exists under this curve in the interval 1 x < .

Ex.) Is the area under the curve y =

3
5

from 1 x < +1 finite? If so, calculate

the value of this area.

In this case, while the limits of integration are finite, there is a point of infinite
1
3
discontinuity at x = 0 . Thus, the integral
dx must be defined as two limits one
5
1 x
that approaches 0 from the left and one that approaches 0 from the right:

5 5 x4
5 5 x4
a
1

A=
dx = lim
dx + lim
dx = 3 lim
1 + 3 lim+
b
5
5
5

4
4
a 0
b0
a 0
b 0
1 x
1 x
b x

4
5 4
5 5 (a ) 4 5 5 (1) 4
5
+ 3 lim 5 (1) 5 (b) =
= 3 lim

4
4
4
4
a 0
b 0 +

5 5
15 15

3 0 + 3 0 = +
= 0 . Note that while the integral equals zero, the area
4 4
4 4

3
15 15
below
is equal to 2 =
. To see why this is the case, observe the graph of y = 5
4
2
x
on the interval 1 x < +1 :

While one of the shaded regions would lie below the x-axis while the other would
lie above it, there is no geometric meaning to a negative area. Thus, the area under this
curve on the given interval is finite, with a value of 15/2.
Applications of Improper Integrals: While it may seem that the significance of improper
integrals is limited because they either involve limits of integration of infinity or points of
infinite discontinuity (or both), this is far from true. Before delving into the physical and
life science applications of improper integrals, it is important to consider them from
somewhat of a philosophical standpoint. This is where the strangeness of calculus comes

20

into play. While our intuition maintains that an infinite geometric figure cannot possibly
have a finite area or volume (after allit is infinite!), the methods of calculus beg to
differ. This is philosophically quite interesting. If one considers the non-Platonistic
premise that mathematics is a tool created by human beings, it is remarkable to note that
the human mind has birthed something that can explain something that it itself cannot
explain! Consider the case of Gabriels Horn, for instance. This geometric figure,
conceived by the Italian mathematician Evangelista Torricelli, has a finite volume, but an
infinite surface area. It can be modeled as the solid of revolution generated when a
1
hyperbola (e.g. y = ) is revolved about the x-axis:
x

Generated on Mathematica

To calculate the volume of this figure, on can use the methods of AP Calculus AB. In
this case, the disc method is the most appropriate:
2

a
1
1
dx
1
1
V =
= lim 1a = lim
+ = (0 + 1) = .
dx = lim
2
a ( a )
a 1 x
a x
(1)
1 x
Thus, the solid has a finite volume. But what of the surface area? While a topic of
neither the AP Calculus AB curriculum nor the AP Calculus BC curriculum, the surface
area of revolution (SA) for a curve f ( x) = y about the x-axis bounded on the closed

interval a x b is described as: SA =

SA =

= 2 lim

dy
2y 1 + dx . For Gabriels Horn:
dx

2
1
1 + 2 dx = 2 lim
a
x
x

1
1+ 4
x
x

dx

1
1
1 1

ln x 2 1a = 2 lim ln a
+ 3 dx = 2 alim

2x
2(a) 2
x x

1
ln 1
2(1) 2

= 2 ( 0) 0 = 2 () = . Thus, the surface area is infinite. This is an

2

interesting scenario. If Gabriels horn were to hold paint, it would hold cubic units of

21

it. Obviously, this is not enough paint to cover its own surface, which is infinite! While
this rather philosophical problem has little practical significance (actually none), the
theory behind this problem plays an integral role (no pun intended) in the sciences in the
form of probability distributions. Note that the following section is not part of the AP
Calculus BC curriculum, but should be read to gain an appreciation of the material.
Probability Distributions in the Sciences: Very often in the physical and life sciences,
one must apply the laws of probability and statistics to solve problems. A probability
density function (often denoted as p(x)) is a mathematical representation that yields the
probability of the occurrence of a random variable x. For a probability distribution on
the interval < x < + , the improper integral is normalized, meaning that it is equal
to unity (or 1):

p( x)dx = 1 . Why is this the case? Remember that p(x) represents a

probability. Loosely speaking, if one considers all of the probabilities associated with a
probability density function, they must sum to 100 percent, or unity. Very often in
nature, the probability density function is a Gaussian function, or normal distribution
1 x 2
1
(a bell-shaped curve), represented as: p( x) =
exp
, where is the
2
2
mean and is the standard deviation. Observe the Gaussian function below:

While there are many variations on the Gaussian function (i.e. some may
skewed to one side, flattened, or stretched), it is imperative to note an important
characteristic of this function; it never equals zero. That is, no matter how far along the
x-axis, the probability of the random variables occurrence is always greater then zero, a
truth that has important implications in the sciences. Note the philosophical connection
to the improper integral problems already discussed; while the improper integral actually
represents a finite number (i.e. 1), the probability density function never touches the xaxis!
One of the most important applications of the probability density function to the
sciences is manifested at the molecular level. The physics applied to macroscopic
objects cannot be used to explain the behavior of the submicroscopic world. Rather, the
constituent components of matter must be described according to the laws of probability
and statistics. For example, consider a collection of gas molecules in a closed container.
If the gas is ideal, then it follows the postulates of the kinetic-molecular theory of gases.
This is a very important concept in the study of physical chemistry, for it describes the
submicroscopic behavior of matter. Where does the probability density function come
22

into play? Within the container described are many gas particles, perhaps on the order of
1023. The kinetic-molecular theory maintains that these particles move randomly in
straight lines and transfer energy and momentum during collisions but do not lose it.
All of the particles are moving at different speeds; one may be moving extremely
quickly, while another extremely slowly and these speeds change quite often.
However, at a given temperature there will be a most probable speed for the collection of
gas molecules. This is shown graphically as a Maxwell-Boltzmann distribution:

Note that the probability never decreases to zero as one moves farther along the x-axis.
Indeed, it is possible for a gas particle at 273 K (0C) to have a speed of, say, 1 10 50 m/s,
but it is not very probable! Notice the characteristics of the Maxwell-Boltzmann
distribution when the temperature increases:

The curve is flattening out such that the most probable speed is greater but is less
probable than the most probable speed at a lower temperature (i.e. the peak is lower).
Why is this case? Recall that for a probability density function the improper integral,
+

p( x)dx = 1 , holds true. Thus, in order to keep the area equal to unity when the most

probable speed is greater in magnitude, the probability of the most probable speed must
be lowered. What would the Maxwell-Boltzmann distribution look like at absolute zero
(0 K)? What about at a very, very high temperature?
Note that the Maxwell-Boltzmann distribution does not solely apply to gases, but
to all matter. This distribution is also important in understanding such processes as
diffusion, which is quite important to life on Earth.
The last application of the probability density function to be considered in this
chapter concerns the field of quantum mechanics. While this is, indeed, a broad field, it
isgenerally based upon the assumption that matter at the submicroscopic level, especially
electrons, exhibit the characteristics of waves more so than particles. As a result of these
wave characteristics, it is impossible to know with certainty both the position and
momentum of a particle simultaneously, a concept known as Heisenbergs uncertainty
23

principle. An electron, for instance, cannot be modeled as a particle, but as a wave

function, denoted as ( x, y, z , t ) , which is a function of three-dimensional space and
time. Squaring this function results in the probability density function. As the distance
from the nucleus increases, the probability of locating an electron greatly decreases, but
never reaches zero. For instance, for the hydrogen atom, its one electron can be
considered a cloud of probable positions in a spherical space:

While it is very unlikely that an electron will be located far away from the
nucleus, it is possible nonetheless. An electron associated with an atom in this sheet of
paper could be somewhere on the moon, but it is not very probable!
Concluding Remarks
This first chapter introduced several new techniques in differential and integral
calculus. LHpitals Rule, integration by parts, integration by partial fractions, and
improper integrals are all very useful tools when faced with more difficult problems that
cannot be solved through the use of the techniques of AP Calculus AB. While these
techniques are powerful, it is important to realize that they are not omnipotent. There are
many other techniques that can be used to evaluate troublesome limits and integrals, and
most are far beyond the scope of AP Calculus BC.
One very important application of improper integrals, the probability density
function, was also discussed in the context of its significance in modeling systems in
physical chemistry and quantum mechanics.
The next chapter will apply many of the techniques discussed in this chapter to
approach various scientific problems.
Key Terms:
indeterminate form
LHpitals Rule
integration by parts
LIPET
tabular method
partial fraction decomposition
repeated factor
integration by partial fractions
reduced row echelon form (RREF)
improper integral
convergent

divergent
Gabriels Horn
probability density function
random variable
normalized
Gaussian function
normal distribution
kinetic-molecular theory
Maxwell-Boltzmann distribution
Heisenberg uncertainty principle
wave function

24

Chapter 2: Differential Equations

Differential equations are introduced in AP Calculus AB. These equations
dy
contain the derivative of a function as a variable. For instance, the equation
= x 2 y is
dx
differential equation that relates the derivative of y with respect to x to those variables
themselves. The differential equations studied in AP Calculus AB and BC are known as
first-order differential equations, because they involve the first derivative of a function.
Differential equations of higher order require more sophisticated techniques to solve and
will not be discussed here, bute will be briefly discussed in the context of chapter 4.
Differential equations are one of the most important (arguably, the most important)
expressions used in mathematical modeling, mostly due to the fact that, by their very
nature, they express a change. This chapter will discuss two ways in which one may
solve differential equations: the method of separation of variables (an analytical method)
and Eulers method (a numerical method). While the AP Calculus AB exam only tests
the first method, the AP Calculus BC exam tests both. In addition to the presentation of
these methods, various applications of differential equations to the physical and life
sciences will be discussed.
As a final note, it is important to appreciate the shear profundity of the study of
differential equations. This chapter only covers two techniques for solving these
equations, which is fine for the AP exam. However, new theories on differential
equations are constantly being developed, and even a course that is specifically devoted
to differential equations will probably not do them justice.
The Method of Separation of Variables
This section should be a review of material from AP Calculus AB. The method
of separation of variables basically separates a differential equation into two sides of an
dy 2 y
=
is a differential equation
equation that have the same variables. For instance,
dx
x
that can be easily solved through algebraically rearranging the equation so that the yvariables appear on one side and the x-variables appear on the other side:
dy 2 y
dy dx
=

= . Once this is completed, one may integrate both sides to determine

dx
x
2y
x
dy
dx
1
ln y
2 ln x
y as a function of x:
=
ln y = ln x e
=e
y = x 2 + C.
2y
x
2
In order to determine the value of C, it would be necessary to be given an initial
condition, such as y (0) = 3 (3) = (0) 2 + C C = 3.

Eulers Method
Eulers Method, pronounced oi-ler, is a numerical method used to approximate
solutions to differential equations. The method was developed by the Swiss
mathematician Leonhard Euler in 1768. Numerical methods for solving differential

25

equations are necessary when solutions cannot be found analytically, such as when one
does not explicitly know the algebraic structure of the differential equation, but knows
certain values for variables and slopes. Given the initial values of the variables and the
slope, a discretized (i.e. non-continuous) form of the limit definition of the derivative can
be used and rearranged to yield the formula for using Eulers method. Recall from AP
y ( x + h) y ( x )
. If h is not
Calculus AB the limit definition of the derivative: y ' ( x) = lim
h 0
h
y ( x + h) y ( x )
infinitely small (i.e. it has a finite value), the equation becomes: y ' ( x)
.
h
For the purposes of Eulers Method, let h be called x , let y ( x + x) be denoted as y n +1
and let y (x) be denoted as y n . The discretized form of the limit definition for the
y n +1 y n
. This is rearranged to yield the formula for
x
Eulers Method: y n +1 = y n + x y ' ( x) n . What is the significance of this formula? If one
knows initial values for a function and its derivative, while not necessarily knowing what
those functions are, and chooses a certain increment x , known as the step size, for the
numerical analysis, one can approximate the value of y(x) for a certain x. See the
example below.
Ex.) Approximate the value of y(1.5) by using increments of 0.1 if y(1) = 4 and
y(x) = x 2 y .

derivative now becomes: y ' ( x)

Notice that in this problem, the actual algebraic structure for the differential
equation is known. This sort of problem often appears on the AP exam, primarily as a
free-response problem. The rationale for a problem of this type will become clear at the
end of this example.
Since the starting x-value is 1 and the step size is 0.1, there must be six steps
involved to approximate a value for the function at x = 1.5. The first conditions given are
x0 = 1 and y 0 = 4 . Thus, the initial slope is (1) 2 (4) = 4 . The formula for Eulers
Method may now be used: y n +1 = y n + x y ' ( x) n y1 = (4) + (0.1)(4) = 4.4.
Five more steps to go! The next x-value is found by adding the step size to the previous
x-value: x1 = x0 + x = (1) + (0.1) = 1.1 . Now y1 can be found:

y1 = x1 y1 = (1.1) 2 (4.4) = 5.324. The remainder of the problem proceeds as follows:

Find y2:
y 2 = y1 + xy1 = (4.4) + (0.1)(5.324) = 4.9324
Find y3:
x 2 = x1 + x = (1.1) + (0.1) = 1.2
2

y 2 = x 2 y 2 = (1.2) 2 (4.9324) = 7.102656

y 3 = y 2 + xy 2 = (4.49324) + (0.1)(7.102656) = 5.2035056
Find y4:
x3 = x 2 + x = (1.2) + (0.1) = 1.3
2

y 3 = x3 y 3 = (1.3) 2 (5.2035056) = 8.793924464

2

26

y 4 = y 3 + xy 3 = (5.2035056) + (0.1)(8.793924464) = 6.082898046

Find y5:
x 4 = x3 + x = (1.3) + (0.1) = 1.4
y 4 = x 4 y 4 = (1.4) 2 (6.082898046) = 11.92248017
y 5 = y 4 + xy 4 = (6.082898046) + (0.1)(11.92248017) = 7.275146063
Find y6:
x5 = x 4 + x = (1.4) + (0.1) = 1.5
2

y 5 = x5 y 5 = (1.5) 2 (7.275146063) = 16.36907864

y 6 = y5 + xy 5 = (7.275146063) + (0.1)(16.36907864) = 8.912053927
Since this last number is the value of y when x = 1.5, this is the approximation for
which the problem asked.
The reason that this sort of problem actually supplies the algebraic structure of the
differential equation is so that one can compare this approximate value with the actual
value yielded from an analytical solution. Indeed, the differential equation given can be
solved analytically:
2

x
x +C
x
3

x3
2

x dx ln y = + C y = e
= e 3 (e C ) = Ce 3 .

C
Recall that C is merely an unknown constant, so e is just another C. This constant
can be determined based upon the initial values given:

dy
= x2 y
dx

y = Ce

x3
3

dy
=
y

(4) = Ce

(1) 3
3

C =

4
e

1/ 3

(1.5 ) 3
3

4
= 8.828287262
y = 1 / 3 e
e
This analytical solution differs from the numerical solution by 0.0837666646.
There are several key points to take away from this Eulers Method problem. Firstly, the
numerical method never yields exact solutions to differential equations. Secondly, as the
step size becomes infinitely small, the solution becomes exact. Thus, the smaller the step
size that one uses, the less error is involved in the calculation. However, notice that even
this problem with the relatively large step size of 0.1 was quite tedious to solve. One
must weigh the precision of the solution needed against the cost of actually solving the
equation numerically. Thirdly, Eulers Method is only one of many numerical methods
of solving first-order differential equations, and is generally the least powerful. In fact,
Eulers Method really only has historical significance, since no one uses this method any
more.
The Law of Exponential Change
Very often in the physical and life sciences, one encounters natural logarithms in
mathematical models of phenomena. This is the case because many aspects of nature
seem to change in direct proportion to an amount of something present. Whether this
amount refers to number of organisms in a population or the number of molecules in a

27

mixture of chemicals in a beaker, the change in these entities often takes on the form of
dN
the following differential equation:
= kN . This mathematical statement means that
dt
the change in some amount N over time is proportional to that amount. That is, if there is
more of something present, its rate of change is greater. To understand why so many
natural processes are modeled mathematically with natural logarithms (in fact, its natural
prevalence is one of the reasons for its name!), it is necessary to solve the differential
equation:
dN
dN
= kN
= kdt ln N = kt + C N = e ( kt +C ) = e kt e C = Ce kt .
N
dt
Let the initial value of N be N0: N 0 = Ce k ( 0) = C .

( )( )

The solution to this differential equation, N = N 0 e kt , is known as the law of

exponential change. The law of exponential change, its derivation, and its applications
are covered in both the AP Calculus AB and BC curricula.
Logistic Growth
While many processes seem to exhibit exponential growth, given enough time,
most of these processes would occur less rapidly once some limiting value is reached.
How can the differential equation for exponential growth be modified to account for this
limiting value? In order to answer this question, it is necessary to translate the meaning
of this logistic growth into a mathematical statement. In the logistic growth model, the
rate of change of a quantity is proportional to both the amount present and the amount
relative to the limiting value, or carrying capacity (K). This is represented by the
dN
following differential equation:
= kN (K N ) . This means that the rate of change of
dt
a quantity N is jointly proportional to the amount present and to the amount relative to the
carrying capacity that is still available for growth. While this differential equation can be
solved by the method of separation of variables, it requires a technique discussed in the
previous chapter integration by partial fractions:
dN
dN
dN
= kN (K N )
= kdt
= kdt = kt +C
dt
N (K N )
N (K N )
1
A
B
= +
Partial fraction decomposition:
N (K N ) N K N
1 = A( K N ) + BN = AK AN + BN = N ( B A) + AK
Since there are no Ns on the left side of the equation, there cannot be any
on the right side either. Therefore, B A = 0. Also, to make the
constants on both sides of the equation equal to 1, A must equal 1/K,
which means that B must equal 1/K as well. Thus:
1
A
B
1
1
= +
=
+
. Integration may now be
N ( K N ) N K N KN K ( K N )

28

carried out:
N
dN
1 1
1
1
kt + C =
=
= Kkt kC
+
dN = (ln N ln K N ) ln
N (K N ) K N K N
K
KN
(As always, kC will just be considered another C)
e Kkt
N
N
K
. This is logistic growth
= kC = Ce Kkt 1 = Ce Kkt N =
KN e
K
1 + Ce Kkt
equation. Note some important characteristics of its algebraic structure. The carrying
capacity is located in the numerator, the exponential is multiplied by a constant C, and e
is raised to the negative power of the product of the carrying capacity, another constant k,
and time. This structure puts a limit on the growth that the equation models. Compare
the exponential growth model to the logistic growth model:

Exponential Growth Model

Logistic Growth Model

Note that for the logistic curve g(t), lim g (t ) = K . The practical applications of
t

these curves will be discussed shortly. Note that logistic growth is frequently tested on
the AP exam either as a multiple-choice or free-response question.
The Learning Curve
In certain cases, the rate of change of a process is directly proportional to the
difference between the endpoint of a process and the degree to which the process has
already ensued. Mathematically, this is expressed as the differential equation:
dN
= k ( A N ), where A is the endpoint of the process. Suppose that the process in
dt
question is the typing of written statement on a word processor. Let the total number of
letters in the statement be 1,000 and let the efficiency be determined by the number of
correct letters typed relative to the 1,000 letters. If the process follows a learning curve,
or bounded growth, the efficiency will increase with time. What is the nature of this
increase? To answer this question, the differential equation must be solved:
dN
dN
dN
= k(A N)
= kdt
= kdt
dt
(A N)
(A N)
ln A N = kt + C ( A N ) = e kt C = (e kt )(e C ) = Ce kt

N = A Ce kt .
In the typewriting example, A would be 1,000 and C and k could be determined
from a set of initial conditions. Notice the graphical nature of the learning curve:

29

The rate of change is most rapid near the beginning and decreases to zero at N=A.
What does this mean? Basically, the greatest degree of learning takes place near the
beginning and decreases as the endpoint is reached. Thus, the more one learns (or the
more of a process that is carried out), the lesser the rate of increase of learning. Note also
that the efficiency, measured in this case as the difference between A and N, increases
with the time invested in the process.
The learning curve has many implications in economic and educational strategy.
If it is the case that as the amount of time invested in a project increases the efficiency
increases, then corporations and classrooms should concentrate on those strategies that
increase the experience of those involved. Problems concerning the learning curve
appear on both the AP Calculus AB and BC exams.
Ex.) An elementary school student must memorize the capitals of all fifty United
States in thirty minutes. The rate of memorization is directly proportional to the
difference between the number of capitals that must be memorized and the number of
capitals that have been memorized so far. Assuming that the child can flawlessly
memorize two state capitals during the first three minutes, and assuming that the child
knew none of the state capitals at the beginning of the study session, will he be able to
memorize all fifty state capitals in thirty minutes?
The problem did not require solving of the differential equation, so it would be
helpful in a problem like this to commit the resulting equation to memory. The same
applies to the exponential and logistic growth models as well, unless, of course, the
problem does ask for the derivation.
N = A Ce kt = 50 Ce kt
It is assumed that the child knew none of the state capitals before beginning
memorization, so (0) = 50 Ce ( 0 ) C = 50 .
ln(48 / 50)
0.0136073.
It is given that at t = 3, N = 2: (2) = 50 50e k ( 3) k =
3
Now, the value of N when t = 30 may be found:
N = 50 50e ( 0.0136073)(30) = 16.75834 17. Unfortunately, thirty minutes is not
enough for the youngster to memorize all of the state capitals.
Mathematical Models of Population Ecology

30

Ecology is a rather broad field of biology that studies the relationships between
organisms and their environment. One branch of this field, population ecology, focuses
on these relationships at the level of a group of the same species inhabiting the same area,
a group referred to as a population. Population ecology is perhaps the most quantitative
of the subfields of ecology, since it often studies the changes in the number of individuals
in a population over time, the strategies that different species use to proliferate, and how
biotic (living) and abiotic (non-living) factors affect a certain population. Populations are
generally modeled mathematically by the exponential growth model or the logistic
growth model.
The Exponential Growth Model in Population Ecology: The exponential growth model
for populations, specifically human populations, was devised by English economist
Thomas Malthus in his An Essay on the Principle of Population (1798). This treatise
exclaimed certain danger for the human race in that while the global food supply grows
linearly, the human population grows exponentially. Population ecologists model this
dN
= rmax N , which is essentially the same equation as the one
exponential growth as
dt
discussed earlier, but instead of using the constant k, rmax is used. This constant is known
as the intrinsic rate of increase, a measure of the capacity of a population to grow at its
maximum potential.
Ex.) A population of fruit flies (genus Drosophila) exhibits exponential growth.
There are initially 10 flies in the population. After three days, the population has
grown to 67 individuals. What is the intrinsic rate of increase of this population?
How many individuals will be present after one week has elapsed?
The problem did not require a derivation of the formula for the law of
exponential change, so it suffices just to commit the formula to memory. For an
example concerning population ecology, the formula is N = N 0 e rmaxt .
It is given that N 0 = 10 and that at t = 3, N = 67. From this information,
the intrinsic rate of increase can be determined:
ln(67 / 10)
(67) = (10)e rmax ( 3) rmax =
0.634036 .
3
Now the value of N when t = 7 can be found:
N = 10e ( 0.634036)( 7 ) = 846.268 846 flies!
At this point, it is probably a good idea to discuss the units in these problems.
The AP Calculus tests are usually not very strict in regards to including units in
calculations, as long as they are provided along with the final answer. Nevertheless,
especially when dealing with functions such as logarithms and exponents, it is important
to discuss the units of each element used in these mathematical models. In the formula,
N = N 0 e rmaxt , N and N0 have units of number of individuals, in the case of the previous
example, flies. The element t has units of time (seconds, minutes, hours, etc.), days in
the previous example. What about rmax? This element is not unitless. In the exponential
growth equation, it has units of reciprocal time, 1/days or days-1 in the previous example.
31

This is because exponentials; along with logarithms, trigonometric functions, and inverse
trigonometric functions; are known as transcendental functions, which are functions
that do not satisfy a polynomial expression. One of the consequences of the failure to
satisfy this expression is that the arguments of transcendental functions must be unitless.
Thus, in the example of the exponential growth model, the product of t and rmax must be
unitless. Since t has units of time, rmax must have units of reciprocal time.
The Logistic Growth Model in Population Ecology: While certain populations exhibit
exponential growth under certain conditions, no population can ever truly attain its
intrinsic rate of increase. A population could only truly grow in an exponential fashion if
it had unlimited access to resources such as food, water, and space. Obviously, these
resources are limited, as Malthus noted, meaning that there is a certain limit to the
amount of individuals in a population that an area can maintain. Indeed, if every
population underwent exponential growth, the total mass of existing organisms would far
exceed the mass of the Earth! Thus, the logistic growth model, discussed earlier in this
chapter, is a far more realistic way to gauge the trends in population growth. Population
dN
(K N )
ecologists often model logistic growth as:
= rmax N
, which is essentially the
dt
K
exponential growth model with another term added to tame the equation. This term,
(K N )
, represents the fraction of the population relative to the carrying capacity that is
K
still available for growth. Note that this is different from the logistic model already
discussed, in which the term was merely (K N), or the number of individuals that can
still undergo growth. While both differential equations yield a logistic function, the one
discussed first should be used for the purposes of AP Calculus BC.
Ex.) A population of grizzly bears in a preserve exhibits logistic growth defined
dN N
N
by the following differential equation:
= 1 . What is the carrying capacity of
dt
6 14
this population? How many grizzly bears are there when the rate of increase in the
population begin to decrease? If when t = 0 years, N = 2 grizzlies and when
t = 5 years, N = 10 grizzlies, determine the time at which this decrease occurs.
For the first part of the problem, recall that the carrying capacity is reached when
the rate of change of the population (the derivative) is equal to zero. Thus, the carrying
capacity can be found directly by setting the differential equation equal to zero are
N
N
solving for N: 1 = 0 N = 14 . This is the carrying capacity for the population.
6 14
Upon reading the second part of the problem, the reader may realize that this
question is essentially testing the ability to analyze the relationships between functions
and their derivatives, an AP Calculus AB skill. When the rate of increase (the derivative)
begins to decrease, the derivative will have a relative maximum. To determine the value
of N at which this occurs one could either graph the derivative on the TI-84 and
determine where its maximum lies, or use the second derivative test. To save time
(assuming that this is a question in which one is permitted to use a calculator!), one
should use the former method. Graphing the derivative, it is found that the derivative has
32

a relative maximum at N = 7. Thus, the rate of change of the population begins to

decrease when there are 7 grizzly bears.
While it may seem that this last problem requires solving the differential equation
dN
(notice that it has a slightly different algebraic structure from
= kN ( K N )) , the
dt
problem stated that the population follows the logistic growth model. So, again, if this
formula is committed to memory, and if the problem does not ask for the derivation, one
can simply use the formula: Two initial conditions are given. These will be used to
K
determine the constants k and C in the equation for logistic growth N =
.
1 + Ce Kkt
14
14
(2) =
=
C =6
14 k ( 0 )
1+ C
1 + Ce
14
ln (4 / 60)
(10) =
14 = 10 + 60e 70 k
= k 0.038686
14 k ( 5 )
70
1 + 6e
ln(7 / 42)
14
14 = 7 + 42e 0.541604t t =
3.30825 .
(7) =
(14 )( 0.038686 ) t
0.541604
1 + 6e
Thus, the rate of change in the population will begin to decrease after a little over three
years.
Mathematical Models of Reaction Kinetics
While the application of exponential and logistic growth will appear on the AP
exam, this following material is not a part of the AP Calculus BC curriculum. However,
it still gives the reader a chance to appreciate a very common application of simple
differential equations.
Reaction kinetics is a branch of physical chemistry that studies the rates and
mechanisms of chemical reactions. This field is of the utmost importance to practical
chemistry because it allows one to determine how quickly a desired reaction will take
place, how the reaction takes place, and how one might speed up (or slow down) the
reaction. Consider the following generic chemical equation:
aA + bB cC + dD

where the capital letters refer to particular chemical species and the lower-case letters
refer to their relative numbers in the reaction (called a stoichiometric coefficient). The
left side of the equation is referred to as the reactants and the right side of the equation is
referred to as the products. As the reaction takes place, the amount of reactants will
decrease and the amount of products will increase. How can one express the rate at which
this reaction takes place? In essence, one can take the derivative of the equation. There
is a problem, however; not all species involved may change at the same rate. For
instance, the coefficient b may be twice as great as coefficient a, meaning B will decrease
twice as fast as A. One can compensate for this by multiplying the rate of change of the
concentration of each species (concentrations are denoted with brackets [ ]) by the
inverse of its stoichiometric coefficient:

33

1 d [ A] 1 d [ B] 1 d [C ] 1 d [ D]

=
+
.
a dt
b dt
c dt
d dt
The negative signs on the left side of the equation signify a decrease in reactants
while the positive sign on the right signifies an increase in products. For instance,
consider the combustion of glucose ( C 6 H 12 O6 ) :
C 6 H 12 O6 + 6O2 6CO2 + 6 H 2 O.
The expression that conveys the reaction rate of this reaction would be:
d [C 6 H 12 O6 ] 1 d [O2 ] 1 d [CO2 ] 1 d [ H 2 O]

=
+
.
dt
6 dt
6 dt
6 dt
The actual nature of reaction rates can be explained by simple differential
equations that can be solved through separation of variables. Reaction rates are usually
gauged by the disappearance of reactants. Depending upon the chemical reaction in
question, these rates are dependant upon reactant concentrations in certain ways. For
instance, in a certain chemical reaction (AB), in which the species A is monitored, the
d [ A]
reaction rate might be expressed by the differential equation:
= k[ A], in which
dt
the rate of decrease of A is proportional to the amount of A present. This equation should
look familiar; its solution is the law of exponential change, but rather than expressing
growth, it expresses decay: [ A] = [ A]0 e kt . A reaction that proceeds in this manner is
said to be of first order because the rate of the reaction is directly proportional to the
concentration of reactants to the first power. In general, a reaction order refers to the
term n in the expression: rate = k [reactants] n , which is known as the rate law of the
reaction. If the rate of decrease of A were proportional to the square of the reactants, the
d [ A]
reaction would follow second-order kinetics (i.e.
= k[ A] 2 ). Reaction rates of
dt
higher order are certainly possible (some biochemical reaction orders are greater than
10), as are fractional orders, zero orders, and negative orders. All of these aspects of the
reaction are determined from laboratory experimentation.

Radiometric Dating
The atom is modeled as a collection of three subatomic particles: protons,
neutrons, and electrons. The electrons are distributed probabilistically (see chapter 1)
around an incredibly dense nucleus, which is composed of protons and neutrons. While
neutral (uncharged) atoms of the same element have the same number of protons and
electrons, they may very well differ in their number of neutrons. Atoms of the same
element that differ in their number of neutrons are referred to as isotopes. The stability
of an isotope is related to the ratio of its number of neutrons to its number of protons. An
unstable isotope is said to be radioactive, and will undergo a mode of radioactive decay
to achieve a more stable neutron-to-proton ratio. Scientists have learned to take
advantage of this submicroscopic phenomenon in the process of radiometric dating, the
determination of the approximate age of an object based upon the amount of a
radioisotope (i.e. radioactive isotope) that it contains, the known decay rate of the
radioisotope, and some reference object that also contains the radioisotope. One of the

34

most revolutionary forms of radiometric dating employs the isotope carbon-14. In this
technique, invented by Nobel-Prize winning chemist Willard Libby in 1949, one
compares the amount of carbon-14 in a sample (measured by the amount of radioactivity)
to the amount in living systems. This is done based upon the premise that the amount of
carbon-14 in carbon dioxide has been relatively constant for many thousands of years
and that the ratio of carbon dioxide with carbon-14 in its molecular structure to carbon
dioxide that does not contain radioactive carbon is the same in the atmosphere as it is in
living things. When a living thing dies, it ceases to assimilate carbon-14 into its
structure, and radioactive decay ensues. Thus, if one measures the amount of carbon-14
in an organically derived sample, such as something made from plant material, and
compares it to the amount of carbon-14 in a living thing, which represents how much
carbon-14 was originally present in the sample, he or she may determine the approximate
amount of time elapsed since the organism whose matter was used to make the object
died. The rate of decay of carbon-14 is not constant, but follows the law of exponential
change. This was determined by the observation that a radioisotopes half-life, the
amount of time elapsed before half of the material has decayed, is independent of any
external factors such as temperature or concentration. The formula for the half-life of an
object can be derived from the law of exponential change:
N = N 0 e kt When half of the material has decayed, half of the material
1
still remains, so N = N 0 . The time at which this occurs is the half-life, denoted
2
1
1
ln 2
as t1 / 2 N 0 = N 0 e kt ln = kt ln(1) ln(2) = kt t1 / 2 =
.
k
2
2
Notice that the half-life only depends upon the constant k, which is a
characteristic only of the specific isotope in question; other factors do not affect
half-life. This holds for other natural systems that follow the law of exponential
change as well. The half-life of carbon-14 is about 5700 years. Note that carbon
dating can only be applied to relatively young samples that are organic in origin.
After about 36,000 years, so much of the original radioactive sample has decayed
that it becomes quite difficult to detect, and the accuracy of the experiment is
substantially reduced. Problems concerning radiometric dating are common on
both the AP Calculus AB and BC exams.
Ex.) A medieval art collector plans to buy a flag that is claimed to have been used
in the Battle of Hastings in 1066. Before paying a hefty price of $50,000 for this
supposed relic, the collector demands that the seller have the flag carbon-dated. Analysis
shows that the radioactivity of carbon-14 detected in living plants from which such a
cloth may be made is 20.9 disintegrations per minute per gram and that the radioactivity
of the material in the flag is 18.6 disintegrations per minute per gram. Note that the halflife of carbon-14 is 5700 years. Is the flag authentic?
Let the variable A represent radioactivity. The radioactive decay of carbon-14
follows the law of exponential change: A = A0 e kt . While the initial radioactivity of the
sample (taken to be equivalent to the radioactivity of the living plant) and the final

35

radioactivity of the sample are known, the constant k is not. This constant may be found
ln 2
ln 2
from the formula for half-life: t1 / 2 =
k=
= 1.216048 10 4 .
k
5700
The approximate time elapsed may now be determined:
4
A = A0 e kt (18.6) = (20.9)e (1.21604810 ) t t = 958.7416 years . Since the Battle of
Hastings occurred 940 years ago, within a certain margin of error, the flag is probably
authentic.
Newtons Law of Cooling
Unlike the complex computational physics of today, the physics of Sir Isaac
Newton was known for its brevity and elegance. One example of this conciseness is
manifest in a simple law of heat transfer known as Newtons Law of Cooling, though the
model applies to heating as well. This model states that the rate of change of an objects
temperature is directly proportional to the difference between the objects temperature
and the temperature of the immediate environment, assuming that this environment is
large enough that it does not experience a substantial change in temperature.
dT
Mathematically,
= k (T Tenv ), where T is the temperature of the object at a certain
dt
time and Tenv is the temperature of the environment, which is assumed to be constant. If
an object is heating up, the value of k will be positive. If an object is cooling down, the
value of k will be negative. Note the important difference between this equation and the
equation for bounded growth; in Newtons Law of Cooling, the endpoint (the
temperature of the environment) is subtracted from the temperature of the object, while in
bounded growth, the variable in question is subtracted from the endpoint. This different
order makes a difference in terms of the final equation, as will be evident in the following
solution to the differential equation for Newtons Law of Cooling:
dT
dT
dT
= k (T Tenv )
= kdt
= kdt
(T Tenv )
(T Tenv )
dt

ln T Tenv = kt + C (T Tenv ) = e kt +C = (e kt )(e C ) = Ce kt .

Note that this is a variation on the law of exponential growth; the only difference
here is that a constant (Tenv) is subtracted from the variable (T). What is the significance
of the constant C? This can be determined through an analysis of initial conditions. At
an initial time t = 0, C = T Tenv . This T is the initial temperature, or T0. Thus, the
equation for Newtons Law of Cooling becomes: (T Tenv ) = (T0 Tenv )e kt . Again, if the
object is heating up, k is found to be positive, and if it is cooling down, it is found to be
negative. Problems concerning Newtons Law of Cooling appear on both the AP
Calculus AB and BC exams.
Ex.) A blacksmith removes a piece of iron from a furnace with a temperature of
266C. The temperature of the blacksmiths workshop is 24C. When the iron has been
allowed to cool for one minute, its temperature is 235C. The blacksmith may work with
the iron when it cools to a temperature of 75C. Assuming that the rate at which the iron
cools is proportional to the difference between the temperature of the iron and the

36

temperature of the environment, how many minutes must the blacksmith wait until he can
handle the iron?
It is given that Tenv = 24 and that T0 = 266. To determine the value of k, one may
use the condition that when t = 1, T = 235 :
211
(T Tenv ) = (T0 Tenv )e kt (235 24) = (266 24)e k (1) k = ln
0.1370796
242
It can now be determined at what time the iron will cool to the desired
temperature:
ln (51 / 242)
11.35918.
0.1370796
Therefore, the blacksmith must wait approximately 11 minutes before handling the iron.
(T Tenv ) = (T0 Tenv )e kt (75 24) = (266 24)e ( 0.1370796 )t

Motion with Air Resistance

According to legend, the famous astronomer Galileo Galilei once conducted an
experiment in which he compared the rate at which two balls fell from the Leaning
Tower of Pisa. One was made of iron and the other was made of wood. Which one
would hit the ground first? If both started at the same height and were released with the
same initial velocity, they should both hit the ground at the same time because they both
experience the same acceleration due to gravity, which is approximately 9.8 m/s2.
Galileo discovered, however, that the iron ball hit the ground before the wooden ball did.
This is because the force of gravity is not the only force acting on these objects; air
resistance also plays a role. However, unlike gravity, air resistance is not a constant
force, but one that is proportional to the velocity of the object. For an object moving at a
relatively slow velocity, such as a feather through the air, the force of air resistance,
called drag (denoted here as Fd ) , is directly proportional to the velocity of the object:
Fd = cv , where c is a constant that depends upon the properties of the air (or another
fluid). An object moving through the air at a relatively fast velocity, such as a speeding
bullet, experiences a drag that is directly proportional to the square of its velocity. While
a thorough discussion of vectors must wait until chapter 6, it is important to introduce
them in the context of this section. Force is a vector quantity, meaning it has both
magnitude and direction. While the choice of a frame of reference is arbitrary, it will be
defined here in an intuitive sense: a downward direction corresponds to a negative vector
and an upward direction corresponds to a positive vector. With this frame of reference,
an object falling through the air is acted upon by gravity (a negative vector) and air
resistance (a positive vector), and is accelerating downward. The following free-body
diagram, which is used in classical physics to depict the forces acting upon an object
(represented here by a dot), should make this clear:

Fd

Fg

37

From this free-body diagram, one can apply Newtons Second Law of Motion,
which states that the net vector sum of the forces acting on an object is equal to that
objects mass times its velocity: F = ma . In the scenario of an object falling through

the air with a relatively slow velocity: F = cv mg = ma , where m is the mass of the

object and g is the acceleration due to gravity, which is approximately 9.8 m/s2 . Recall
from AP Calculus AB that acceleration is the derivative of velocity. Thus, Newtons
dv
Second Law of Motion is essentially a differential equation: ma = m
= cv mg .
dt
While this differential equation is not tested on the AP exam, it still makes for a good
exercise in solving differential equations by the method of separation of variables while
also providing a meaningful application. This differential equation may be solved as
follows:
dv
mdv
mdv
m
= mg cv
= dt
= dt.
dt
mg cv
mg cv
These indefinite integrals can be made into definite integrals to simplify the
calculations. Assume that the initial velocity is zero (at time t = 0) and that any velocity
at time t is v. Then:
v
t
mdv
= dt = t t0 = t.
0 mg cv
0
The integral on the left can be solved through u-substitution. Let u = mg cv.
c
Then, du = -cdv, and the integrand must be multiplied by in order to yield the form
m
m
u / du = ln u + C , meaning the outside of the integral must be multiplied by c . Now:
m
m
v
t = ln mg cv 0 = [ln(mg cv) ln(mg )]
c
c

(mg cv)
mg cv
cv
mg
ct
= ln
e ct / m =
= 1
v=
(1 e ct / m ).
m
(mg )
mg
mg
c

This is a function of velocity with respect to time for the object as it falls through
the air. Notice the characteristics of this function when it is graphed:

38

The rate of change of velocity (the acceleration) decreases until a constant

velocity is reached. This constant velocity is known as the terminal velocity.
Note that drag may also be proportional to the square of the objects
velocity if this velocity is of relatively large magnitude. The differential equation that
results from this scenario requires integration that yields a certain kind of function (an
inverse hyperbolic tangent, or tanh 1 ) . While it is not terribly difficult to solve, it is
somewhat beyond the scope of this book.
Concluding Remarks
This chapter can be considered a brief introduction to differential equations and
their power as mathematical models of natural phenomena. Note that while this chapter
only discussed first-order differential equations and two methods for solving them (the
method of separation of variables and Eulers Method), there are many more analytical
and numerical approaches.
Hopefully this chapter also instilled in the reader a greater appreciation for the
applications of differential equations to the physical and life sciences.
Key Terms:
first-order differential equations
method of separation of variables
Eulers Method
step size
Law of Exponential Change
logistic growth
carrying capacity
learning curve
bounded growth
population ecology
intrinsic rate of increase
transcendental functions
reaction kinetics
reaction order
rate law

isotopes
radioactive decay
radiometric dating
half-life
Newtons Law of Cooling
drag
free-body diagram
Newtons Second Law of Motion
terminal velocity

39

Chapter 3: Infinite Sequences and Series

Thus far, this book has been dedicated to derivatives and integrals, whether it be
using LHpitals Rule, integrating using new techniques, or solving differential
equations. For many students, the next two chapters contain material that is quite
puzzling and quite boring puzzling because it seems to relate very little to other topics
in AP Calculus AB and BC, and boring because, most of the time, it introduces
absolutely no applications. In this book, I will try to remedy these problems by
conveying both the relation of the material to the primary theme of differential and
integral calculus, moving to the level of the infinitesimal, and its most important
applications.
Infinite Sequences
Unlike functions, which have been studied thus far, sequences only include
certain elements defined by a particular rule that is performed upon positive integers. In
this respect, sequences are discretized, whereas functions take on many values on a
certain interval. A sequence is basically a list of integers that have been affected in the
same way, meaning that they follow a general pattern. A sequence a1 , a 2 , a3 , L, a n may
be written more succinctly as {a n }. When studying a sequence, one must know the rule
that generated the terms of the sequence. For instance the first five terms of the sequence
n 2 , in which n is a positive integer, are 1, 4, 9, 16, 25, L . In the study of calculus, one is
concerned with the behavior of infinite sequences, which contain an infinite number of
terms. Some series will approach a certain value (i.e. they will have a limit), whereas
others will continue their pattern of growth without approaching any particular value.
Those sequences that approach a certain value are convergent, while those that do not are
divergent. While the limits of sequences can be evaluated in much the same way as the
limits of functions, it is important to discuss a theorem on infinite sequences and how
they specifically relate to functions in regards to determining limits. This theorem has
three postulates:
1.) If there exists an infinite sequence {a n } and a function f(x) such that
for every integer n, f (n) = a n , and
2.) if the function f(x) is defined for all x 1 , and
3.) if lim f ( x) exists, then lim a n must exist, and lim f ( x) = lim a n .

{ }

While one cannot evaluate limits of infinite sequences in the same manner as
functions, one can relate a certain infinite sequence to a function through this theorem,
and by determining the limit of the function, one can determine the limit of the sequence.
ln n
Ex.) Does the infinite sequence 3 have a limit as n approaches infinity? If
n
so, what is this limit?
The question is asking if the sequence converges or diverges. If the limit
exists, the sequence converges. If the limit does not exist, it diverges. In order to
40

determine this limit, it is necessary to employ the theorem that relates infinite sequences
ln x ln()
= . This limit is in indeterminate form. Therefore,
to functions: lim 3 =
3
x x

( )
ln x
1/ x
1 1 1
= = 0.
LHpitals Rule must be used: lim 3 = lim 2 = lim 2 =
2
x x

x 3 x x 3 x 3()
ln n
Thus, lim 3 = 0.
n n

This answer could also have been determined through the use of the TI-84. If one
presses the MODE key and chooses seq, he or she may graph sequences rather than
functions. Going to Y=, inserting 1 as nMin, typing ln(n)/n3, in the u(n) row,
and pressing the GRAPH key (one might need to choose an appropriate window for
viewing), one should notice two key features of the resulting graph. First, it seems to
approach the x-axis, which accords with the previous analytical solution of 0. Second,
the graph is not an unbroken curve, but a collection of dots. Since this is the graph of the
sequence, only certain, discrete values can be assumed, which are represented by the
dots:

Infinite Series
While sequences are merely lists of numbers, series are sums of the terms of a
sequence. In essence, infinite series and sequences are perhaps even more closely related
than this statement implies. Consider the sequence a1 , a 2 , a3 , a 4 ,L, a n . Suppose one
were to increasingly sum the terms of this sequence, i.e., start with a1 , then add a1 to a 2 ,
then add a1 , a 2 , and a3 , and so on. Each of these instances of increasingly encompassing
summation is called a partial sum, because, indeed, it is only a partial sum of the terms
of an infinite sequence. Thus, a new sequence is formed: S1 , S 2 , S 3 ,L, S n , in which
S1 = a1 , S 2 = a1 + a 2 , S 3 = a1 + a 2 + a3 , and S n = a1 + a 2 + a3 + L + a n . This new
sequence is an infinite series, an infinite sum of the terms of an infinite sequence. Series
are often not indicated as a sequence of partial sums, but in sigma notation:

= a1 + a 2 + a3 + L + a n .

n =1

41

At this point, a certain aspect of infinite series arises that defies intuition, but
which, nonetheless has a very real basis in calculus. While infinite series are the sum of
an infinite number of terms, this sum can be a finite number. Recall the discussion of
Zenos paradoxes from the prelude. The theory behind infinite series can actually solve
many of these paradoxes quite elegantly, as will become evident soon in this chapter. An
infinite series whose terms sum to a finite value is said to be convergent, while an infinite
series whose terms sum to infinity (i.e. a finite sum does not exist) is said to be divergent.

There are two ways in which an infinite series may converge. If a series

a is
n

n =1

absolutely convergent, then the sum of the absolute values of the terms is also

convergent, i.e.,

is convergent. If the series

n =1

a converges, but a
n

n =1

does

n =1

not, the series

a is said to be conditionally convergent. This distinction will

n

n =1

become more meaningful as the chapter progresses, particularly because while the
convergence of a conditionally convergent series depends upon the order in which the
terms are summed, this makes no difference for an absolutely convergent series. Lastly,
there are three helpful algebraic properties of infinite series that are exactly analogous to
the algebraic properties of integrals:

1.)

(a + b ) = a + b
n

n =1

2.)

(a

n =1

n =1

3.)

bn ) =

n =1

a b
n

n =1

n =1

ca = c a .
n

n =1

n =1

Determination of Infinite Series Convergence or Divergence

There are many techniques used to determine the convergence or divergence of a
series, each of which partly depend upon the nature of series. While these techniques
will easily convey whether a series has a finite sum, most of them are not powerful
enough to determine what the sum actually is. Note that examples of problems using
convergence and divergence tests will appear after the full discussion of these tests.
Nth Term Divergence Test: The nth term divergence test is one of the simplest
tests to use when one begins his or her inspection of a series. It must be used with
caution, however. This test only proves that a series diverges; it cannot prove

convergence. This test states that if for the series

a , lim a
n

n =1

0, then the series

diverges. However, if lim a n = 0, this does not prove convergence; it means that another
n

42

test must be used. It is easy to understand this from a logical standpoint; if the statement
If the limit does not equal zero, then the series diverges, is true, the reverse statement
If the limit equals zero, then the series converges, is not necessarily true. Only the
contrapositive of the original statement is definitely true: If the series converges, then
the limit equals zero. However, this statement cannot be used in the nth term divergence
test because the flow of logic must move from the truth about the limit to the conclusion;
the logic of the contrapositive statement moves from the conclusion to the truth about the
limit. If one knew the conclusion beforehand, one would already have the answer!
As a precaution, it is important to note that the limits that are evaluated in these
tests do not represent the actual sums of the series; their values merely convey whether
the series in question converges or diverges.
Test for Geometric Series: A geometric series is a series in which all terms have
a common factor and a ratio r that is raised to a specific power. Mathematically, a

geometric series has the form

ar

= a + ar + ar 2 + L ar n . Notice that the lower-

n =0

bound limit in the summation symbol is 0. This is often the case so that the common
factor a can be found easily, since it would be the first term in the sequence. There is
another way to make for this simplicity; a geometric series may sometimes be written as

ar

n 1

= a + ar + ar 2 + L ar n . This also allows for the first term to be a. These

n =1

reconstructions bring up an important point; a removal or addition of terms (as long as

that removal or addition is finite!) will not affect the convergence or divergence of an
infinite series. A geometric series will converge if 0 < r < 1. While the proof will not
be shown here, geometric series are one of the very few series whose exact sum can be
found. The formula for this infinite sum (denoted as S n , the nth partial sum), provided the
a
geometric series converges, is S n =
. Infinite geometric series are quite important in
1 r
the sciences, as many discrete systems go through geometric progression. In addition, an
infinite geometric series can be used to solve one of Zenos paradoxes, in which the
philosopher muses that motion is futile, for in order for an individual to move a length l,
he or she must move half that distance, and half that distance, and half that distance, and
so on. In essence, Zenos paradox is an infinite geometric series. The geometric series
represents the various distances that the individual must travel 1/2, 1/4, 1/16, 1/32, etc.
One can find the value of a and r by applying the equation for an infinite geometric
series: since a is the first term, a will equal 1/2 in this case. The second term is ar, which
is 1/4 in this case. Thus, r must also equal 1/2. The infinite geometric series for Zenos

n
1
1 1
1 1 1 1 1
+L+
. Since
paradox is thus:
= + + +
2 2n
2 2 2 4 16 32
n =0
1
1
r = 0 < < 1 , the series does have a finite sum. Using the formula for the nth
2
2

43

a
(1 / 2)
=
= 1. Thus,
1 r 1 (1 / 2)
motion is not futile as the paradox would purport because the infinite number of distances
actually sum to a whole!
Test for P Series: A p-series is an infinite series in which the terms to be added
partial sum to find what this finite sum actually is: S n =

together are risen to a negative exponent:

n1

n =1

1
1
1
1
+ p + p + L + p . A p-series
p
1
2
3
n

will converge if p > 1 , and will diverge if p 1. An interesting case that lies right on the
border of these conditions is the harmonic series, a p-series with p = 1, and, thus,
divergent. This name is derived from the physics of a vibrating string, which vibrates at
integer multiples of a fundamental frequency. These discrete multiples follow a p-series
in which p = 1. This is the mathematical basis for harmony in music. When one plucks
or bows a string instrument, or plays a wind instrument for that matter, it is not the case
that only one frequency is generated; most music would sound awful if this were how
acoustics worked! Instead, one actually perceives an infinite sum of frequencies that
follow a harmonic series (mostly some instruments do deviate slightly from the
harmonic series).
Telescoping Series: Besides the geometric series, there is one other series whose
exact sum can be very easily found the telescoping series, in which all but one or a few
terms will cancel each other through subtraction. Telescoping series are given their
informal name based upon the fact that they collapse like a pocket telescope, meaning
that most of the terms of this series disappear, leaving only a remnant of the fully
expanded form of the series. Telescoping series contain an infinite number of terms and
their opposites, so one can often spot such series by the presence of subtracted terms. See
the example concerning telescoping series at the end of this discussion of tests for
convergence and divergence.
The Integral Test: Recall from the beginning of this chapter that it is not possible
to perform the same calculus on discrete expressions like sequences and series as on
functions. However, there are mathematical theorems that allow one to analyze
discretized expressions as if they were functions. One such theorem is the integral test,

which states that if the series

a is a monotonic decreasing series, meaning that it

n

n =1

continually decreases, with positive terms, and that f(x) is also a monotonic decreasing,
positive function, which is also continuous, such that f (n) = a n for every positive integer

n, then the series

a n converges or diverges if the improper integral

n =1

f ( x)dx
1

converges or diverges, respectively. The proof will not be shown here.

The Ratio Test: For series whose terms contain factorials (!), it is often useful to

employ the ratio test. The theorem of this test states that a series

a converges when
n

n =1

44

a n +1
a
< 1 and diverges when lim n +1 > 1. However, if the limit is equal to 1, the test
n a
n a
n
n
is inconclusive. This test is particularly helpful in areas such as computer science and
combinatorics, a branch of discrete mathematics that studies (as its name suggests)
combinations of entities. This is because the factorial is an operation that expresses the
different ways in which entities may be arranged.
The Root Test: While the root test is rarely considered on the AP Calculus BC
exam, it is quite powerful. It is very closely related to the ratio test and is often
considered more powerful than the ratio test, because where the ratio test fails the root
lim

test often succeeds. The theorem of the root test states that a series

a converges if
n

n =1

lim n a n < 1 and diverges when lim n a n > 1. Similar to the ratio test, if the limit is
n

equal to 1, the test is inconclusive. Thus, since this makes the series appear smaller
than it would by using the ratio test, a series that appears too large to be convergent in
the ratio test might actually be proved convergent in the root test.
The Comparison Test: If one considers two infinite series, one which is known to
be convergent or divergent and one which is not, he or she can determine if the
unknown series converges or diverges by comparing it to the known series. If there

are two series

a and b that have positive terms, and every term in b is

n

n =1

n =1

n =1

greater than the corresponding term in

a , then the theorem for the comparison test

n

n =1

holds true:

a converges if b converges (i.e. the smaller series will definitely

n

n =1

n =1

converge if the larger series converges, but not necessarily vice versa), and

b diverges if a diverges (i.e. the larger series will definitely diverge if

n

n =1

n =1

the smaller series diverges, but not necessarily vice versa).

This should make intuitive sense: If a series that has larger terms converges to
some value, it is certainly the case that a series with smaller terms must also converge to
some value. Furthermore, if a series with smaller terms does not approach a value, then a
series with larger terms definitely will not approach a value. Note that in order to carry
the comparison test out validly, it must be known with certainty that for every positive
integer n, one series has every corresponding term greater than another series. For

instance, consider the two series

(2
n =1

+ 1 and

(2

+ 3 . Obviously, for every n,

n =1

the second series will have larger terms. This might not be so clear on other occasions.
The Limit Comparison Test: Oftentimes, one is not completely sure that every
term in one series is greater than the corresponding term in another series. However, the
45

series can still be compared through the limit comparison test. The theorem of this test

states that if there are two series

a and b that have positive terms, such that

n

n =1

n =1

an
= L , where L 0 and L , then both series behave in the same way, i.e. they
n b
n
either both converge or both diverge.
The Alternating Series Test: Besides the telescoping series, which could have
negative terms, all of the series discussed so far have had positive terms. In an
alternating series, the terms oscillate between positive and negative values, making
these series suitable for modeling oscillatory phenomena. Alternating series are
characterized by the presence of 1 raised to some power. For instance,
lim

(1) a , (1)
n

n =1

n 1

a n , or

n =1

(1)

n +1

a n . In order for an alternating series to

n =1

converge, three conditions must be fulfilled: all a n s are positive, a n > a n +1 for all n, and
lim a n = 0 . Do not become confused by the first part of the theorem; while the a n s
n

cannot be negative, the terms most certainly can, and will be! The negative 1 to an
exponent is not part of the a n .
It is now appropriate to bring up the topic of absolute and conditional

convergence again. Recall that a series is absolutely convergent if

is convergent

n =1

and conditionally convergent if

a is convergent but not a

n

n =1

. If a series has been

n =1

proved to be absolutely convergent, then it is also convergent in general. Technically,

if

a converges but not a

n

n =1

, then the series really cannot be referred to as

n =1

convergent; it must be specified as conditionally convergent. Many of the tests

discussed prove absolute convergence immediately. This is because these tests contain
an absolute value in the formula for determination. These include the tests for geometric
series, the ratio test, and the root test. For the other tests, if a question asks for
specification of absolute or conditional convergence, the absolute value must be taken as
an extra step.

Ex.) Does the series

56n

converge or diverge?

n =1

In the analysis of series, it is always helpful to algebraically clean up

a n as best as possible. In this problem, 6/5 can be put in front of the summation symbol:

46

6
5

n1 . Once this is done, one should first attempt to use the n

th

term

n =1

divergence test; if this test proves that the series diverges, one need go no further.
Remember, however, that this test does not prove convergence. Applying this test on this
1
1
series, lim 3 =
= 0. Thus, the nth term divergence test is inconclusive; another test
2
n n
( )
is needed. This series is easily recognizable as a p-series with p = 3. Since 3>1, the series
converges.
Note that, with enough mathematical experience, one could recognize this
series as a p-series and skip the nth term divergence test altogether. However, the next
problem will show that it can be quite useful for more algebraically complicated series.

Ex.) Does the series

n =1

3n 2 + 6n 1
converge or diverge?
n2

While more sophisticated tests could be used, one would benefit from first
th
applying the n term divergence to a complicated a n such as this:
3n 2 + 6n 1
lim
= 3. Thus, since this number does not equal zero, the
n
n2
series diverges. This limit was determined through the use of a theorem regarding limits
approaching infinity from AP Calculus AB. Recall the theorem that states that a limit
approaching infinity of a ratio of algebraic statements that have the same order is equal to
the ratio of the leading coefficients (i.e. the numbers multiplied by the variable raised to
the highest power) in the numerator and denominator, in this case, 3 and 1.
9 27 81
+
+ + L. Does this
2 4
8
geometric series converge or diverge? If it converges, what is its sum?
Ex.) The following is a geometric series: 2 + 3 +

It is given that the sum of the terms represents a geometric series. In a

geometric series, the first term is always a, so a = 2 in this case. Furthermore, the
3
second term is ar. Since ar = 3 in this case, and since a = 2 , r = . Thus this
2

n
3
3
geometric series can be represented as
2 . Since > 1, the series diverges.
2
2
n =0

Ex.) Does the series

4 +14n

converge or diverge?

n =1

Trying the nth term divergence test results in a limit of 0. Thus, another
test may be used. This is a power series with an extra term added in the denominator.
While other tests could be used, the comparison test will yield an answer more quickly.

47

Let

n =1

1
be the smaller series and let
4 + 4n 2

n =1

1
be the larger series. One can
4n 2

easily determine if the larger series diverges because it is a simple p-series. Since p = 2,
which is greater than 1, the larger series converges. Thus, according to the theorem for

the comparison test, the smaller series,

4 +14n

, must converge as well.

n =1

Ex.) Determine the sum of the telescoping series

1n n +1 4 .
n =1

There are no special formulae to determine the answer; the first several
terms must be written out and one must decide which terms cancel and which remain:

1n n +1 4 = 1 15 + 12 16 + 13 17 + 14 18 + 15 19 + 16 101 . Note the

n =1

1 1
, , and, if the series were to be written out in its entirety, an infinite
5 6
1 1
1
number of fractions after that! The only terms that do not cancel are 1, , , and .
2 3
4
1 1 1 25
Thus, the sum of this telescoping series is: 1 + + + = .
2 3 4 12

terms that cancel:

Ex.) Does the series

n =1

n
converge or diverge? If it converges, does it
2e n

converge conditionally or absolutely?

The nth term divergence test would be inconclusive since the limit turns out to be
0. Thus, another test must be used. This is neither a geometric series nor a power series.
The best test to use in this case is the integral test. If one graphs the function
x
f ( x) = x , he or she would confirm that for all positive integers, it is continuous,
2e
monotonic decreasing, and positive. Thus, the integral test can be used. One creates the

x
improper integral
dx and rearranges it algebraically for ease of analysis:
x
1 2e

a
1 x
1
xe dx = lim
xe x dx. This integral must be solved through the method of
2 1
2 a 1
integration by parts. One notices that one of the terms is a polynomial (x), so the tabular
method can be used:

48

u
x

1
lim
2 a

xe x dx =

dv

1
+1

ex

-1

ex

+1
-1

1
lim xe x e x
a
2

a
1

{[

][

]}

1
12 1
lim (a )e ( a ) e ( a ) (1)e (1) e (1) = = .
2 a
2e e

Thus, since the integral converges, the series also converges.

The next part of the question asked for a specification of the kind of convergence.

It must be determined if

n =1

n
converges or diverges. For this, one could use the
2e n

comparison test, in which the known series for comparison is

n =1

and greater than

n =1

n
, which is divergent
1

n
for all positive n. This comparison series is divergent by the
2e n

nth term divergence test. Thus, the comparison test is inconclusive and one should move
on to the limit comparison test:
n
1
1
1
2e n
= lim n =
=
= 0. Thus, the series do not behave in
lim
(

)
n
n 2e
n

2e
1

the same way, so the original series,

n =1

n
, must converge. Therefore,
2e n

n =1

n
is
2e n

absolutely convergent.

Ex.) Does the series

10n!

4n

converge or diverge?

n =1

Since this series contains a factorial, the ratio test is probably the best
choice here:
(n + 1)!
4( n+1)
4n
10
= lim (n + 1)! 10 Note the property of factorials that
lim
4( n+1)

n!
n
n 10
n!

4n
10

(n + 1)! 10 4 n
10 4 n (n + 1)n!
(n + 1)
=
= lim
(n + 1)!= (n + 1)n!. Thus, lim 4( n +1)
. The
.
lim

n
+
4
(
1
)
n 10
n (10 4 )
n! n 10
n!

algebra was a bit cumbersome to yield this last fraction. While the n!s cancelled well,

( )

the other terms required an application of the following laws of exponents: a bn = a b

49

ab
(n + 1)
= a b c . That aside, lim
= . Since infinity is greater than 1, the series
c
n (10 4 )
a
diverges.

and

Ex.) Does the series

n
2

n =1

+1

converge or diverge?

The nth term divergence test is of no use here because the limit is equal to
zero. This is neither a geometric series nor a p-series. While other tests could be used on
this series, the least cumbersome is probably the limit comparison test. This series can be

compared to the simpler series

nn , which is, in fact, a p-series with p = 32 . Thus,

2

n =1

is not equal to zero or

n2
n2 +1
infinity, the series act in the same way; in this case, they would both converge:
n

n2 +1
n2 n

n5/ 2

. Since this is a
lim
= lim 2
= lim 5 / 2
1/ 2
n n
n n
n n
n
n
n
+
+

n2

limit involving infinity in which the orders (i.e. highest powers) of the numerator and
denominator are the same, one takes the ratio of the leading coefficients, which is 1 in
this case. Thus, the limit is exists and is non-zero. Therefore, both series behave in the

this series converges. If the limit of the ratio of

same way and

n
n =1

n
2

+1

to

converges.

n 2 2n 1
2
converge or diverge?
Ex.) Does the series
+

5
n
16
n
12

n =1
This series would be easy to analyze if the n power were not present. A
good rule-of-thumb is, if the series has at least one term raised to the n with no factorials
present, the root test is probably the best bet. This is the test that will be used in this
problem:
n

n 2 2n 1
n 2 2n 1
= lim 2
lim n 2
. Notice how this gets rid of
n
n 5n + 16n 12
5n + 16n 12
the troublesome n exponent. Since this is now a limit involving infinity with the same
order in the numerator and denominator, one can take the ratio of the leading coefficients:
n 2 2n 1
1
1
lim 2
= . Since < 1, the series converges.
n 5n + 16n 12
5
5

50

Ex.) Does the series

n =1

n
converge or diverge?
(5) n 1

This is, in fact, an alternating series that requires some algebraic

manipulation to become easier to work with:

n 1

n
1
(1) n 1 n 1 .
n =
5
5
n =1
n =1
n =1
Now one can decide whether this series is convergent or divergent by
using the three criteria described for an alternating series:
Every a n is positive.
n
=
(5) n 1

?
?
5 n ? (n + 1)
(n + 1)
n
n 1

(
5
)(
n
)
>
(
5
)(
n
+
1
)

>

5
n
>
(n + 1)
n
5 n 1
5 n 1
5n
?
? 1
1
4n > 1 n > . Yes, if n 1 , every n will be greater than .
4
4
( )
n

lim a n = 0 lim n 1 = ( ) 1 = Use LHpitals Rule:

n
n 5

5
1
1
1
n
lim n 1 = lim n 1
= ( ) 1
= = 0.
n 5
n 5
ln 5 5
ln 5

a n > a n +1

>

All three criteria are fulfilled, so the series converges.

Approximations of the Sums of Infinite Series
Determining if a series converges or diverges is a very important skill in
mathematical analysis and computational science. However, the next step can be slightly
more challenging actually determining what the sum is. While there are a number of
techniques that one may employ to determine the sum of an infinite series, all of them are
approximations. Nevertheless, for the purposes for which one uses these series, either in
pure or applied mathematics, these approximations are suitable. It is even theoretically
impossible to sum an infinite series to completion because infinity is not really a number.
This, yet again, introduces the troubling concept of the infinitesimal that is ubiquitous in
calculus. Thus, even when exact sums of geometric and telescoping series were found,
these were still limits. This section will introduce the approximation methods tested on
the AP Calculus BC exam and will discuss a very consequential application of infinite
series approximations to an important problem in mathematics. Note that there are a
myriad of other techniques used to approximate the sums of infinite series, many of
which are carried out via computer.

51

Approximation by Truncation: Consider the geometric series

n =0

a
of this series can be determined analytically: S n =
=
1 r

1
2 . The exact sum
2

= 4. How many terms of

1
1
2
this series would one have to add in order to approach a reasonable approximation? This
is the question asked when dealing with infinite series whose exact sum cannot be found.
In truth, this depends upon the requirements of the problem; in some cases extreme
precision is necessary, while in others, a few summed terms will suffice. The
approximation referred to here is known as truncation, which, as the name suggests, is a
separation of a certain number of terms of an infinite series from the rest of the terms.
These chosen terms are then summed to yield an approximation of the value to which the

1
2 such that the first five terms
series converges. Consider truncating the series
2
n =0
are summed, which is referred to as the fifth partial sum. How close is this value to the
1 1 1
actual value of 4? S 5 = 2 + 1 + + + = 3.875. While this seems rather close to the
2 4 8
actual value, it is helpful to define a parameter that will unequivocally indicate how far
off the estimate is. This parameter is known as the truncation error, which is, of course,
the error associated with splitting up the infinite series and summing a representative
number of terms. One can define the truncation in this context (Note that the term
truncation error will take on a slightly different meaning in the next chapter) to be the
percent error between the exact sum and the estimated sum. For the fifth partial sum:
4 3.875
error =
100 = 3.125%. Relatively speaking, this is a good estimate. While
4
this example conveyed the theory behind truncation, it was rather meaningless in a
practical sense since the exact sum of the geometric series was known. It is far more
interesting to consider examples in which only approximation works, but in which one
can also determine the truncation error.
The case of practical truncations to be discussed here is often tested on the AP
Calculus BC exam, and applies to alternating series. While the proof will not be shown
here, a theorem of alternating series states that the absolute value of the maximum

possible truncation error for a series

a , in which the sum is approximated by

n

n =1

summing the first i terms, is less than or equal to the value of the (i + 1)st term: Ri ai +1 ,
where Ri is the truncation error associated with taking the sum of the first i terms.

Ex.) Calculate the error bound when the sum of the alternating series
th

n =1

(1) n n
10 n

is approximated with the 100 partial sum.

52

While the problem implies that this alternating series converges, one must
always first prove that it does indeed converge using the alternating series test:
Every a n is positive.
n ? (n + 1)
>
Since an exponential function
10 n
10 n +1
increases much more rapidly than a linear function, the denominators of these
fractions will become larger more quickly than the numerators. Since the
n (n + 1)
(n + 1)
magnitude of the exponential for
is greater than that for n ,
will
n +1
10
10
10 n +1
be smaller than the other fraction because of its rapidly increasing exponential of
n
(n + 1)
higher magnitude in the denominator. Thus, n >
.
10
10 n +1
Prove that a n +1 < a n :

Prove that lim a n = 0 : lim

n

n 10 n

= 0 true.

Thus, the series does indeed converge.

If the sum of this series is approximated with the 100th partial sum, the
maximum possible error must be less than or equal to the value of the 101st term:
(102)
R101 a102 R101 (102 ) . Obviously, this error is extremely small so
10
small that the TI-84 does not know what to make of it! Thus, summing the first 100
terms of this alternating series is a very good approximation of the exact sum.

Ex.) What is the sum of the alternating series

(1)
n =1

n 1

n
5 n 1

to an accuracy of 5

decimal places?
This is the series from the example in the last section, which was proved to
converge. Thus, the proof need not be carried out here.
The problem asks for an accuracy of 5 decimal places, which means that
Ri 0.00001. From this information, one must determine how many terms (i) must be
summed to approximate the value within 5 decimal places:
j
Ri ai +1 a i +1 = 0.00001 Let i + 1 = j j 0.00001 . This is quite
5
difficult to solve analytically, because the variable is tied up both in a linear function and
in an exponential function. It is a good idea to solve this graphically using the TI-84.
x
One graphs y = x and y = 0.00001 and determines where they intersect. These two
5
graphs intersect at x = 8.4817448. Thus, (i + 1) must be less than or equal to 8.4817448.
The value i must be the next integer less than 8.4817448, which is 7. Therefore, by
summing the first 7 terms of the series, one achieves an accuracy within 5 decimal places:

53

(1)
n =1

n 1

n 1

(1) 0

7
6
5
4
3
2
1
+ (1)1 1 + (1) 2 2 + (1) 3 3 + (1) 4 4 + (1) 5 5 + (1) 6
0
5
5
5
5
5
5
5

1 0.4 + 0.12 0.032 + 0.008 0.00192 + 0.00044 = 0.69452.

A Case Study of Infinite Series Approximation: Will We Ever Know The Value of ? :
It is quite possibly the most enigmatic numerical expression of all human history.
For millennia, mathematicians and scientists have arduously striven to determine the
value of to more and more decimal places. Most of the numerical techniques
employed in these toils involve infinite series. While the Egyptians and Babylonians
were quite successful in their computations, extending the value of to several decimal
places by 1000 CE, it was arguably not until the late 14th century when the venerable
Indian mathematician Madhava of Sangamagrama, controversially considered by some
to be the true founder of calculus, developed a highly effective infinite series with which
he determined to 11 decimal places:

1
(1) n 1 12

1 1
= 12 1 +

+ L = . During the early 20th century,

n 1
(2n 1)(3 )

9 45 189
n =1
another Indian mathematician, Srinivasa Ramanujan, developed another infinite series
that has spearheaded the modern computational advances in the calculation of .
Ramanujans infinite series, though enormously helpful to other mathematicians, is not

very elegant:

2 (1103 + 26390n)
29801
[(n!) (396) ] .
4

4n

Contemporary technology has extended the

n =0

known value of by millions of decimal places. Perhaps the continually increasing

prowess in humankinds computational ability will one day lead to the discovery of
something upon which philosophers and mathematicians have ruminated since the
discovery of this elusive constant: Is truly an irrational number, or does it actually
terminate at some point beyond the horizon?
Concluding Remarks
This chapter introduced the concept and theory of infinite sequences and series.
The majority of the chapter stressed the various tests for convergence and divergence of
infinite series, in which situations these tests are appropriate, and in which situations they
are not.
While these tests allow one to determine whether or not a series has a finite sum,
they cannot immediately convey what that sum actually is, and only sums for geometric
and telescoping series can be evaluated exactly in an analytical sense. All is not lost,
however. There are many suitable methods one can use to approximate the sum of an
infinite series with great accuracy and precision. While only two were discussed in this
chapter, the method of truncation and the theorem of alternating series, there are a great
number of very powerful computational methods that form the basis of the study of
numerical analysis. The evolution of these techniques throughout the course of human
54

history are manifested in the continual computational search for more and more digits to
.

Key Terms:
sequences
infinite sequences
series
partial sum
infinite series
absolutely convergent
conditionally convergent
nth term divergence test
geometric series
p-series
harmonic series
telescoping series
integral test
ratio test
combinatorics
root test
comparison test
limit comparison test
alternating series

truncation
truncation error
Madhava of Sangamagrama
Srinivasa Ramanujan

55

Chapter 4: Power Series and Polynomial

Approximations
The material in this chapter is extremely important, not just for the purposes of a
thorough preparation for the AP Calculus BC exam, but for an understanding of some
extremely useful mathematical tools used in nearly every subfield of applied mathematics
and the sciences. The majority of this chapter emphasizes the importance and power of
approximations in quantitative approaches to understanding phenomena. One can often
lose sight of the usefulness of such approximations when he or she wields such a
sophisticated device as the TI-84 that can spit out strings of decimal places. Interestingly,
though, the algorithms of calculators use some of the same techniques that are discussed
in this chapter! One must not dismiss the importance of these techniques, as is
unfortunately so often the case when students are unaware of their consequential
applications.
Power Series
The previous chapter introduced infinite sequences and series of constant terms.
This chapter will concern infinite series of non-constant, variable terms. Such series are

known as power series, and are generally denoted as

a ( x c) , where a is an
n

n =0

expression with ns, the same as the expressions associated with the series discussed in the
previous chapter, x is a variable, and c is a constant at which the series is centered. Note
that power series have lower limits of 0. Besides the fact that power series contain
variable terms (i.e. xs), there is yet another intricacy that must be introduced: It does not
suffice to specify that a power series converges; one must also specify where it
converges. A power series may converges in three ways: 1.) It could converge only at the
value c at which it is centered, 2.) It could converge at all real numbers within a certain
radius r (called the radius of convergence) from the center c (and, further, this could
include the value r of the radius, or it could not), or 3.) It could converge at all real
numbers. The complexity of the second case must be discussed further. Suppose the
graph below represents a power series centered at the origin (i.e. c = 0):

The radius to the right of c has a value of (c + r ) and the radius to the left of c has
a value of (c r ) . Note that c is positioned at the origin here for simplification; a power

56

series could have a non-zero value of c as well. If the power series represented in the
graph converges within the radius of convergence, it must be specified if this does or
does not include the values of (c + r ) and (c r ) . That is, it must be specified whether
the interval of convergence, all of the values at which the series converges, is a closed
interval or an open interval, an intricacy discussed in Calculus AB or, for those students
who have taken the class, pre-calculus. Note that while the term radius of convergence
is usually only applied to situations like that implied in the graph, it could have a value of
zero if it is only convergent at c or it could have a value of infinity if it is convergent at
all real numbers.
How can one determine the radius of convergence of a power series? In general,
one uses the ratio test introduced in the previous chapter because it gives concrete
a
numerical conditions for convergence; the series will converge when lim n +1 < 1 .
n a
n
Ex.) Determine the radius and interval of convergence of the power series

n =0

n 2 ( x 1) n
.
4n

Recall that power series have the form

a ( x c) . In this is example,
n

n =0

n
and c is equal to 1. One can use the ratio test to determine at
4n
what values of x this series converges:

the a n term is

(n + 1) 2 ( x 1) ( n +1)

4 ( n +1)
lim
n
n 2 ( x 1) n

4n

2
( n +1)
4n
= lim (n + 1) ( x 1)
n
4 ( n +1)
n 2 ( x 1) n

(n + 1) 2 ( x 1) (n 2 + 2n + 1)( x 1)
= lim
=
4n 2
4n 2
n

1
= ( x 1). Note that the (x 1) remains in the same form because it does not
4
contain any ns, so it is not part of the limit. The ratio test maintains that a series
a
1
converges when lim n +1 < 1 , so the series will converge when ( x 1) < 1 . From this
n a
4
n
expression, one can determine the radius and interval of convergence:
1
1 < ( x 1) < 1 4 < ( x 1) < 4 3 < x < 5 . While the radius of
4
convergence is definitely 4, one must test both endpoints of -3 and 5 to determine if the
interval of convergence includes these points, one must substitute them for x in the
original series:

n =0

n 2 (3 1) n
=
4n

n =0

n 2 (4) n
=
4n

diverges by nth term divergence test

n =0

57

n =0

n 2 (5 1)
=
4n

n =0

n 2 (4) n
=
4n

diverges by nth term divergence test.

n =0

Thus, the interval does not include the endpoints, so the interval of convergence is
solely 3 < x < 5. Note that, in this case, the nth term divergence test was sufficient to
prove divergence when x was -3 and 5. However, for other problems the other tests
discussed in the previous chapter may be necessary.

Ex.) Given the power series

n =0

n p xn
, determine the value, or values, of p such
3n

that the interval of convergence includes the endpoints of the interval.

np
and c = 0. It is given that the radius of
3n
convergence must equal 2. Using the ratio test,
(n + 1) p x ( n +1)

p ( n +1)
3( n +1)
x(n + 1) p 1
3n

= lim (n + 1) x

=
= x
lim
lim
n
n
3
n p x n n 3n p
3( n +1)
n p xn

n
3
In this power series, a n =

1
x < 1 3 < x < 3.
3
One must now determine for what value, of values, of p the following
series will converge:

n =0

n p (3) n
=
3n

n =0

and

n =0

n p (3) n
=
3n

n =0

Both of the above series can be considered the reciprocals of p-series. A

p-series converges for all p greater than 1, so the above series must converge for all p less
than -1.
Power Series as Functions
Thus far, infinite series have been discussed as if they were mathematically
separate from functions. Many power series, however, are actually mathematically
equivalent to functions, that is, expressions that do not contain the ns of a series
1
expression. Consider the very simple example of f ( x) =
. This can actually be
1 x

58

represented as a geometric power series. Recall that the nth partial sum of a geometric
a
1
series is equal to
. Thus, for the function f ( x) =
, a = 1 and r = x. Thus, the
1 r
1 x
function may be expressed as follows:
1
f ( x) =
=
1 x

= 1 + x + x2 + x3 + L xn .

n =0

More complicated series require differentiation and integration to determine their

functional counterparts in terms of x without ns. Consider a function f(x) that is defined

by the generic power series

a ( x c) . There is a theorem that states that if this

n

n =0

power series converges on a certain interval of convergence, the function that it defines
has an infinite number of derivatives inside this interval. If one differentiates each term
of the series, the resulting series will also converge within the same interval of
convergence, but not necessarily including the endpoints. The action of taking the
derivative of each term of such a power series is known as term-by-term
differentiation. This is easier to understand mathematically:

If f ( x) =

a n ( x c) , then f ( x) =
n

n =0

f ( x) =

na ( x c)
n

n 1

n =0

n(n 1)a ( x c)

n2

, etc. converge within the same interval

n =0

of convergence as the original power series.

There is also a theorem that states a similar truth for integration. That is, if a

function f(x) is defined by the power series

a ( x c)

, then the integral of this power

n =0

series,

n =0

a n ( x c)
n +1

n +1

, converges within the same interval convergence as the original

series, but not necessarily at the endpoints. The action of integrating each term of such a
power series is known as term-by-term integration. It is important to note a key
difference between the term-by-term integration theorem and the term-by-term
differentiation theorem; the term-by-term integration theorem does not state that one can
expect an endless number of integrals of a power series to converge within the same
interval convergence as that original power series only the first integral.
How are these theorems useful? Suppose one must know the power series
associated with ln(1 x) but only knows the geometric power series that defines the

x = 1+ x + x + x +Lx .
1
One could integrate this power series term-by-term because
1 xdx = ln(1 x) + C.
1
, which, as derived previously, is
function f ( x) =
1 x

n =0

However, one wishes to cause the constant of integration to disappear in deriving the

59

series for ln(1 x) . Thus, one can perform a definite integral and define f(x) in terms of a
new variable t:
ln(1 x) =

x
0
3

1
dt =
1 t

n =0

x n+1
t2 t3 t4
= t + + L
n +1
2 3 4

x2 x
x4
= x +

+
L
2
3
4
This last expression is the power series that represents ln(1 x). Thus, it is very
helpful to recognize a known power series as a derivative or integral of an unknown
power series. But how are these known power series determined in the first place?
This was a relatively simple task in regards to the geometric series because the function
a
was in the form of S n =
. How might one go about finding an explicit function for
1 r

(1) n x 2 n +1
? In order to simplify this problem, one wishes to
the power series
n
+
2
1
n =0

somehow create a geometric power series so that a function may be easily generated in
a
accordance with the formula S n =
. This is first carried out through term-by-term
1 r
differentiation:

f ( x) =

n =0

(1) n 1 (2n 1) x 2 n 2
=
2n 1

(1)

n 1

x 2 n 2 = 1 x 2 + x 4 x 6 + L.

n =0

Notice that this is merely an alternating geometric series, in which a = 1 and

r = x . Thus, the derivative of f (x) may be expressed in terms of x:
1
1
f ( x) =
=
.
2
1 ( x ) 1 + x 2
To find the original function f (x) , one takes the indefinite integral of its
derivative:
1
f ( x) = f ( x) =
dx = tan 1 x + C , in which the constant of
2
1+ x

(1) n x 2 n +1
is f ( x) = tan 1 x.
integration is 0. Thus, the explicit function for
2n + 1
n =0
2

Unfortunately, this technique cannot be used in every circumstance. That is, it is

not always possible to simplify a complicated series into a more recognizable geometric
series. There are other analytical and graphical techniques used in these circumstances,
and the result is substantial tables of the functions corresponding to specific power series.
For the purposes of the AP Calculus BC exam, the following power series and their
corresponding functions should be memorized:

60

Power Series

((21)n +x1)!
n

2 n +1

x
x
x
+

+L
3
!
5
!
7
!
n =0
Interval of Convergence: < x <
= x

x2 x4 x6
(1) n x 2 n
= 1
+

+L
(2n)!
2! 4! 6!
n =0
Interval of Convergence: < x <

n =0

x3 x5 x7
(1) n x 2 n +1
= x
+

+L
2n + 1
3
5
7

xn
x2 x3 x4
= 1+ x +
+
+
+L
n
!
2
!
3
!
4
!
n =0
Interval of convergence: < x <

n =0

(1) n ( x 1) n +1
( x 1) 2 ( x 1) 3 ( x 1) 4
= ( x 1)
+

+L
n +1
2
3
4

n =0

(1) n x n = 1 x + x 2 x 3 + x 4 + L

ex

ln x

1
1+ x

Interval of Convergence: 1 < x < 1

n =0

tan 1 x

Interval of Convergence: 0 < x 2

cos x

Interval of Convergence: 1 x 1

Function
sin x

(1) n ( x 1) n = 1 ( x 1) + ( x 1) 2 ( x 1) 3 + ( x 1) 4 + L

1
x

Interval of Convergence: 0 < x < 2

An often tested concept on the AP Calculus BC exam is that a certain power series can be
manipulated algebraically to yield another power series. For instance, consider the last
1
power series in the table, for f ( x) = . One can easily find the power series for
x
1
f ( x) = 2 by simply substituting x 2 for every x. Thus,
x

1
=
(1) n ( x 2 1) n = 1 ( x 2 1) + ( x 2 1) 2 ( x 2 1) 3 + ( x 2 1) 4 + L. One can also
2
x
n =0

multiply and divide power series with similar ease.

Maclaurin Series, Taylor Series, and Polynomial Approximations
The conclusion of the last section discussed the methods used to derive an explicit
function from a power series. Here, the reverse approach will be covered. The technique
61

of expressing a function as a power series is one of the most important tools in

computational science. Through the methods discussed in this section, it is possible to
express a complicated, sometimes unapproachable, function as a simple polynomial.
Maclaurin Series: If a certain function f (x) has infinitely many derivatives in the
vicinity of a certain point, the origin, for instance, it can be expressed in the general form
of a power series as follows:

f ( x) =

n =0

x n f ( n ) (0)
x2
x3
= f (0) + xf (0) +
f (0) +
f (0) + L , in which the
n!
2!
3!

(0) indicates that the series is centered about the origin and the (n) denotes the nth
derivative of f (x). Such a power series is known as a Maclaurin series. While such
series had earlier been described by Madhava of Sangamagrama (see previous chapter)
and 17th century mathematician James Gregory, they are named for a later contributor,
the 18th-century Scottish mathematician Colin Maclaurin, who mentioned their use in his
Treatise on Fluxions (1742).
The above power series is influential because, in essence, it simplifies
complicated functions to polynomials of certain order; the higher the order of the
polynomial, the better the approximation to the complicated function in question.
Consider f ( x) = sin x as a simple instance of employing such series. As one increases
the order of the Maclaurin polynomial, the approximation to sin x becomes increasingly
better
1st - order Maclaurin polynomial : x cos(0) = x
x2
2 nd - order Maclaurin polynomial : x cos(0) sin(0) = x
22!
x
x3
x3
3 rd - order Maclaurin polynomial : x cos(0) sin(0) cos(0) = x
2!2
3!3
3!
4
x
x
x
x3
4 th - order Maclaurin polynomial : x cos(0) sin(0) cos(0) +
sin(0) = x
2!
3!
4!
3!
x2
x3
x4
x5
x3 x5
th
+
5 - order Maclaurin polynomial : x cos(0) sin(0) cos(0) + sin(0) + cos(0) = x
2!
3!
4!
5!
3! 5!
Notice in the graph below that as the degree of polynomial increases, so does its
nearness to sin x in the vicinity of the origin:

62

The dark, shaded curve is y = sin x , the solid line is y = x, the dotted curve is
x3
x3 x5
+ . It should be clear that as the degree
, and the solid curve is y = x
3!
3! 5!
of the polynomial increases, so does its agreement with the function that it approximates
in the vicinity of where the Maclaurin series is centered (i.e. x = 0 ).
When solving problems concerning Maclaurin series, it is helpful to make a table
with one column for n, one column for f ( n ) ( x), and one column for f ( n ) (0) to ease the
burden of constructing the Maclaurin polynomial. For instance, if one were to use a 3rdorder Maclaurin polynomial to approximate the value of sin x at 0.01, he or she would
create the following table:
n
f ( n ) ( x)
f ( n ) (0)
sin x
0
0
cos x
1
1
sin x
2
0

cos
x
3
-1
y = x

Using the table as a guide, one can easily generate the Maclaurin polynomial:

n =0

f ( n ) (0) x n (0) x 0 (1) x 1 (0) x 2 (1) x 3

x3
=
+
+
+
= x
n!
0!
1!
2!
3!
3!

(0.01) 3
= 0.0099998333.
3!
Ex.) Using a 3rd-order Maclaurin series approximate the value of tan(0.2) .

Using this polynomial, one approximates sin(0.01) as (0.01)

n
0
1
2

f ( n ) ( x)
tan x
sec 2 x
2(sec 2 x)(tan x)

2 (sec 4 x) + (sec 2 x) + (tan 2 x)

n =0

f ( n ) (0)
0
1
0
4

f ( n ) (0) x n (0) x 0 (1) x 1 (0) x 2 (4) x 3

4x3
2x3
=
+
+
+
= x
= x
n!
0!
1!
2!
3!
3!
3

2(0.2) 3
= 0.1946666667.
3
As a graphical supplement, observe the proximity to tan x of the polynomial in
the vicinity of x = 0 :
tan(0.2) (0.2)

63

The solid curve is tan x , and the dotted curve is the polynomial.
Taylor Series: When Colin Maclaurin published the series that are now named after him,
he knew very well of the toils of an English mathematician named Brook Taylor. In fact,
Maclaurin series are merely special cases of Taylor series; while the former are centered
at the origin, the latter can be centered anywhere. Taylor series are thus used when
approximating values of functions that are not in the vicinity of x = 0. For instance,
while Maclaurin series are valid for x-values such as 0.01 and 0.2, they are not valid for a
value such as x = 2. In such cases, one approximates the complicated function in
question with a Taylor polynomial, which is generated by the following rule:

( x c ) n f ( n ) (c )
( x c) 2 f (c) ( x c) 3 f (c)

= f (c ) + ( x c ) f (c ) +
+
+ L.
n!
2!
3!
n =0

Note that this is essentially the same formula as that of the Maclaurin series,
except that this formula is centered on a non-zero value c.
Ex.) Approximate the value of ln(1.6) using a 5th-degree Taylor polynomial
centered at x = 2.
Just as with the setting up of Maclaurin polynomials, a table greatly facilitates the
creation of Taylor polynomials:
n
0
1
2
3
4

f ( n ) ( x)
ln x
1/ x
1/ x 2
2 / x3
6 / x4

f ( n ) (2)
ln(2)
1/ 2
1/ 4
1/ 4
3/8

24 / x 5

3/ 4

( x 2) 0 ln 2 ( x 2)(1 / 2) ( x 2) 2 (1 / 4) ( x 2) 3 (1 / 4) ( x 2) 4 (3 / 8) ( x 2) 5 (3 / 4)

+
+
+
+
+
0!
1!
2!
3!
4!
5!
2
3
4
5
( x 2) ( x 2)
( x 2)
( x 2)
( x 2)
= ln 2 +

+
.
2
8
24
64
160

64

(1.6 2) (1.6 2) 2 (1.6 2) 3 (1.6 2) 4 (1.6 2) 5

+
= 0.47001651.
2
8
24
64
160
Notice the close proximity of y = ln x and the Taylor polynomial in the vicinity of

ln(1.6) ln(2) +
x = 2:

The solid curve is y = ln x and the dotted curve is the Taylor polynomial.
In the previous chapter, the term truncation was introduced in the context
of splitting an infinite series into a certain number of partial sums while leaving
the rest of the series unaccounted for. In performing such an action, one can
determine the error associated with the chosen partial sum. Truncation and the
error associated with it can be interpreted similarly for Taylor series. However,
one must understand truncation in the context of a very important theorem. In
one of Taylors most influential papers, Methodus incrementorum directa et
inversa (1715), he describes what is now known as Taylors theorem. This
theorem concerns the relationship between the chosen Taylor polynomial, denoted
as Pn (x), and the rest of the infinite series, the remainder, denoted as Rn (x).
Taylors theorem states that if a function f (x) has n derivatives on the closed
interval a x b and its (n + 1) st derivative exists on the open interval a < x < b ,
there exists some number (the Greek letter xi-pronounced zye) in the closed
f ( n +1) ( )( x c) n +1
and
(n + 1)!
f ( x) = Pn ( x) + Rn ( x). Now for the English version: Taylors theorem states that
a function f (x) that may be defined as a power series consists of the chosen
polynomial, Pn (x), and the remainder Rn (x), which makes intuitive sense; there is
a component of the series that one is using and a component that one is not using.
The remainder term Rn (x) is nonetheless quite useful because it allows one to
determine the error associated with the polynomial approximation. It is important
to understand the significance of the terms in the formula for the
f ( n +1) ( )( x c) n +1
remainder: Rn ( x) =
. The term f ( n +1) refers to the (n + 1) st
(n + 1)!
derivative of the function that one is approximating. For instance, if one were
using a 4th-order Taylor polynomial, he or she would evaluate the 5th derivative of
that function for the remainder formula. Note that this derivative is evaluated at a

interval a x b such that Rn ( x) =

65

number , which some books may refer to as z. This number occurs between the
center of the series, c, and the value that one is evaluating, x. In order to
determine the maximum possible error associated with a polynomial
approximation, it is necessary to choose a value of between x and c such that it
makes the error, Rn (x), as large as possible. The formula
f ( n +1) ( )( x c) n +1
represents the Lagrange form of the remainder,
(n + 1)!
named after Joseph Louis Lagrange, an 18th-century Italian-French
mathematician. There is also another form known as the Cauchy form of the
remainder, named after the 19th-century French mathematician Augustin Louis
Cauchy, which involves integration, but will not be discussed here.
Notice the similarity between this error analysis and that which was
introduced in the previous chapter concerning alternating series. Both involve the
use of the (n + 1) st term to determine the truncation error. In the context of Taylor
series, the truncation error, the absolute value of the remainder, must be less than
or equal to the Lagrange term. The following two examples should clarify the use
of the Lagrange form of the remainder to estimate.
Rn ( x ) =

Ex.) Use a 4th-order Taylor polynomial centered at x = 1 to determine the

value of e ( 0.8) . Calculate the Lagrange error bound associated with this
polynomial approximation.
As usual, one creates a table:
f ( n ) ( x)

n
0
1
2
3
4

ex
ex
ex
ex
ex

f ( n ) (1)
E
E
E
E
E

e( x 1) 0 e( x 1)1 e( x 1) 2 e( x 1) 3 e( x 1) 4

+
+
+
+
0!
1!
2!
3!
4!
2
3
4
e( x 1)
e( x 1)
e( x 1)
= e + e( x 1) +
+
+
2!
3!
4!
e(0.8 1) 2 e(0.8 1) 3 e(0.8 1) 4
e e
e
e
e
e + e(0.8 1) +
+
+
=e +

+
=
2!
3!
4!
5 50 750 15000
2.225547942.
( 0.8 )

To determine the truncation error associated with this approximation, the

f ( n +1) ( )( x c) n +1
Lagrange form of the remainder is used: Rn ( x) =
. Since a 4th-order
(n + 1)!
66

( )(0.8 1) 5
. What value of
5!
must be chosen? In order to make the error as large as it can possibly be, one chooses
the largest value of f(x) on the closed interval 0.8 x 1 , which is 1. Thus,
f ( 5) (1)(0.2) 5
e ( 0.2) 5
Max error Rn (0.8) =
=
7.24875154 10 6.
5!
5!
Taylor polynomial was used, (n+1) = 5, and Rn (0.8) =

(5)

In another type of problem, one must determine the order of the Taylor
polynomial to be used based upon the maximum acceptable error.
Ex.) What minimum order Taylor polynomial centered at x = / 2 must be used
to approximate the value of cos( / 5) with an error no greater than 0.0003?
In this problem, c = / 2 and x = / 5. Since n is not known, neither is
st
the (n + 1) derivative known. Thus, the appropriate value of must be chosen
somewhat more hypothetically than in the previous problem. Another complexity is that
differentiation of cos x makes for an oscillation between positive and negative values of
cos x and sin x. In any event, what is the greatest value that a pure (i.e. without
multiplying, dividing, adding, or subtracting constants) sine or cosine function may have?
The answer 1. Which function (i.e. cos x, cos x, sin x, or sin x) has a value of 1
somewhere on the closed interval / 5 x / 2 ? The only one is sin x , and this occurs
at x = / 2. Thus, the (n + 1) st derivative is sin x, and the value for is / 2. However,
the value of n is still not known, but this can be easily found from the formula for the
f ( n +1) ( )( x c) n +1
Lagrange form of the remainder: Rn ( x) =
:
(n + 1)!
Max error Rn ( x) =

f ( n +1) ( )( x c) n +1
(n + 1)!

sin( / 2)( / 5 / 2) n +1
(0.0003) Rn ( / 5) =
(n + 1)!
(3 / 10) n +1
. This would be relatively simple to solve if it were
(n + 1)!
not for the factorial, so one must use trial and error. By testing values of n, one finds that
(3 / 10) ( 6+1)
(3 / 10) (5+1)
n must equal 6, because
= 0.0001311 and
= 0.0009734 .
(6 + 1)!
(5 + 1)!
Thus, one must use a 6th-order Taylor polynomial.
(0.0003)

An important technicality must be introduced at this point. Thus far, it has been
assumed that the Taylor polynomials used all converge within the area in which the
approximation is taking place. For functions such as cos x, sin x, e x , ln x , and other
common functions, this is true. In general, however, convergence is not necessarily

67

guaranteed. However, one can determine if a Taylor series converges in the neighborhood
in which one is approximating by using the formula for the Lagrange form of the
remainder. As the order of the Taylor polynomial increases, the remainder should
become smaller. This should make sense when considering that f ( x) = Pn ( x) + Rn ( x); as
n increases, the term Pn (x) will become larger, which means that the term Rn (x) must
become smaller. In fact, for a convergent Taylor series, when n = , Rn ( x) = 0. This
should also make sense; if an infinite-order Taylor polynomial is used to approximate,
there should be no remainder, and no error! Thus, if lim Rn ( x) = 0, the Taylor series is
n

convergent in the vicinity of c, and perhaps everywhere. This intricacy, though important
theoretically, is not an aspect of the AP Calculus BC curriculum.

Applications of Taylor Series

It seems to be the case that students do not appreciate the significance of the
material in this chapter when it is first presented. For one, why would anyone wish to
approximate the value of cos( / 5) when he or she could merely punch a few keys on the
TI-84? Secondly, if this material is supposedly so consequential, why does the AP
Calculus BC curriculum not cover the applications of Taylor series as it does for other
topics such as differentiation and integration? When I first encountered the material in
this chapter, I also viewed it as quite meaningless and, worse, a form of circular logic; if
we are using our calculators to approximate a value of a complicated function, why
couldnt we just use the calculator to yield a more precise value in the first place? Since
that time, I have researched the importance of Taylor series in applied mathematics and
science and now see that they are among the most important mathematical tools a
scientist can use. Generally speaking, describing a mathematical model as a Taylor
polynomial greatly reduces the complexity of the problem at hand because it is simply
easier to work with polynomials than other functions. In fact, it is often the case that the
use of Taylor polynomials is the only way to approach a mathematical problem. While
this concluding section of the chapter cannot possibly do justice to the vast number of
Taylor series applications, the representative applications discussed here should give the
reader an idea of their power in mathematical modeling.
Simplification to Simple Harmonic Motion: Physical systems such as swinging pendula
and masses attached to springs can generally be described by the same equations. These
equations describe a state of simple harmonic motion, motion characterized by a linear
restoring force, or a force that varies linearly with distance and restores some object to
its original position. In this section, the example of a swinging pendulum will be
considered. It can be shown through dynamical analysis (i.e. analysis of the forces acting
on the pendulum) that the pendulum obeys the following equation:
d 2
g
d 2
=

in
which
sin

,
is the angular acceleration of the
l
dt 2
dt 2
pendulum, g is the acceleration due to gravity, l is the length of the pendulum, and is

68

the angular position of the pendulum. Unfortunately, this equation does not represent
simple harmonic motion, which would state that the second derivative of position (i.e.
acceleration) is equal to a negative constant times the position, which would have the
d 2
g
form of
= . It is the case, however, that the swinging of a pendulum can be
2
l
dt
approximated as simple harmonic motion for small angles. Why is this the case? For
small angles, sin . Thus, as long as the pendulum does not stray too far from the
d 2
g
vertical, it can be described by the equation
= . To understand this
2
l
dt
graphically, consider the fact that the line y = x approximates the value of y = sin x very
well be in the neighborhood of the origin:

d 2
g
This is an important because, relatively speaking, the expression
= is a
2
l
dt
second-order differential equation that can be easily solved while the expression
d 2
g
= sin is basically intractable. In essence, one is performing a polynomial
2
l
dt
approximation to simplify the differential equation at hand. There are many other
situations in which polynomial approximations are necessary in the physical sciences, as
will shown in the next two subsections.
Modeling Molecular Vibrations: All objects in the universe exert forces on one
another. This is true for matter at the submicroscopic level as well. Molecules are
composed of atoms participating in chemical bonds, and these atoms exert both repulsive
and attractive forces upon other nearby atoms. Experiments in the field of quantum
chemistry by the English mathematician John Lennard-Jones in 1931 have led to a
function that conveys how the potential energy of the atoms in a noble gas (e.g. helium
and argon) varies with increasing distance between these atoms. This function, known as
12 6
the Lennard-Jones potential, is given as V (r ) = 4 , where V (r ) is the
r
r
potential energy as a function of the distance between the atoms in a molecule and and
are both constants. Below is a plot of the Lennard-Jones potential, with distance r on
the x-axis and potential V on the y-axis:

69

The Lennard-Jones potential, as one can see from the unwieldy formula, is rather
difficult to work with in computational studies. Thus, one chooses to simplify it as a
Taylor polynomial that approximates the potential energy values in the vicinity of the
distance associated with the lowest potential energy value. The Taylor polynomial
expansion of the Lennard-Jones potential about is:

n =0

(r ) n V ( n ) ( )
(r ) 2
(r ) 3
= V ( ) + V ( )(r ) + V ( )
+ V ( )
+L
n!
2!
3!

Noting that is the minimum x-value, this is the value at which the first
derivative of the Lennard-Jones potential point should be zero. This allows one to find
in terms of :
6 6
V ( ) = 4
7
( )

12 12

13
( )

= 0 6 6 13 = 12 12 7 6 13 = 2 12 7

6
7
=
= 26 .
12
13

2
Thus, the first two terms of the Taylor series are
12 6
V ( ) = 4 6 6 = and zero, since the first derivative at was
2
2

found to be so. To approximate the Lennard-Jones potential by a quadratic Taylor

polynomial, one more term in the series is needed, so one takes the second derivative of
the Lennard-Jones potential at :
156 12
V ( ) = 4 6
14
(2 )

42 6 36 3 4

(2 6 ) 8 = 2 .

36 (r ) 2
.
Thus, the Taylor polynomial is V (r ) +
3
2 2
Observe how closely, even at just two terms, this polynomial approximates the
Lennard-Jones potential in the vicinity of :

70

The thicker curve is the Lennard-Jones potential and the thinner curve is the
second-order Taylor polynomial.
Given the relatively cumbersome calculations involved here, one might ask,
How did finding a Taylor polynomial help us at all? The answer lies in a concept
known as computational cost. While the initial calculations might have seemed
substantial, the resulting polynomial eases subsequent calculations, either by hand or by
computer, to such a degree, that the net payoff is considerable.

The Virial Equation: To the readers who have studied chemistry and physics, the
equation PV = nRT must look familiar. This is known as the ideal gas law, in which P is
pressure, V is volume, n is the number of moles of the gas (a measure of the number of
particles), R is a constant, and T is temperature. As the name of this equation suggests, it
is an idealization that only works in rather limited circumstances. Many great scientists
have modified the ideal gas law to describe the behavior of real gases, but even these
equations are often insufficient. One need not approach the problem with knowledge of
the explicit algebraic structure of the mathematical model, however. In fact, one of the
powerful attributes of Taylor series is that they can function as mathematical models even
when one does not know anything about the actual structure of the equation at hand. In
the case of a real gas, the pressure (P) of the gas in question can be described as a
function of the concentration of the gas particles (M) and temperature (T): P = f ( M , T ).
Notice that this equation is not very specific at all. In fact, it merely states that there is
some function that expresses the relationship between the three variables. This is of little
concern. Assuming that the function is continuous, one may expand it as a Taylor series
1 2 f (M , T ) 2
f ( M , T )
as follows: P = f ( M , T ) +
M+
M +L
2! M 2
M
Do not be put off by the new notation; the symbol means that one is taking what is
known as a partial derivative, because the function in question is one of three variables,
not just two. A (much) more extensive discussion of partial derivatives is reserved for
multivariable calculus. In any event, it is important to note important aspects of the
Taylor polynomial above. First of all, the partial derivatives are merely coefficients of
the variable term M. Secondly, one can evaluate this polynomial at a gas concentration
of zero, which means that M = 0, technically turning the Taylor polynomial into a
Maclaurin polynomial. When the gas concentration is zero, the gas becomes ideal,
n
meaning that the first term, f ( M , T ) , is RT , which can be written as MRT , because M
V
f ( M , T )
2 f (M , T )
just refers to the moles of gas over volume. Since the terms
and
M
M 2
71

are merely coefficients, one can simplify the polynomial by dividing both sides by MRT
and giving the partial derivatives simpler names of C and B:
P
= 1 + BM + CM 2 + L.
MRT
The equation above is known as the virial equation, and B and C are virial coefficients,
which can be determined experimentally or theoretically. One need not know the science
behind this equation to appreciate its significance; almost nothing was known before the
polynomial expansion, but now this equation yields very useful information about the gas
in question, more useful, in fact, than could any of the other algebraically explicit
equations!
In these last two subsections of this chapter, the application of Taylor series to
problems that often arise in mathematics will be discussed.

Taylor Series and Integration: It is often the case that one will encounter integrals that
simply cannot be expressed in terms of familiar functions. Not only are these integrals
difficult to evaluate, but it is impossible to evaluate them in terms of algebraic or
transcendental functions. Consider the following definite integral:

sin( x )dx. There

2

is actually no familiar function whose derivative is sin( x ). (If the reader really must
know, this kind of integral is called a Fresnel integral, which appears often in wave
optics). However, integration of polynomials is the easiest integration there is! The
definite integral above can be approximated as the definite integral of a Maclaurin
polynomial (because the values over which one is integrating are close to the origin).

Recall that the polynomial expansion for sin x is

n =0

x3 x5 x7
(1) n x 2 n +1
= x
+

+ L.
(2n + 1)!
3! 5! 7!

One can simply substitute x for x , allowing one to approximate the definite integral as
follows:
1
1

x 6 x10 x14
sin( x 2 )dx = x 2
+

+ Ldx .

3!
5!
7!
0
0

Similar tricks can be carried out for other difficult integrals.

Taylor Series and Differential Equations: The very broad topic of differential equations
was introduced in chapter 2. There are countless analytical and numerical techniques that
are used to solve differential equations. In chapter 2, the method of separation of
variables and Eulers method was discussed. This subsection introduces another
technique, known as the method of power series, which is quite helpful in solving
second-order differential equations. This method is best explained through a simple
dy
example. Consider the differential equation
= ky. From chapter 2, it was determined
dx
that the solution is an exponential function. While this differential equation can be easily
solved by the method of separation of variables, it is used as an example here to provide a

72

dy
= ky
dx
by the method of power series. To do this, one would first perform a polynomial
expansion upon y as follows: y = a 0 + a1 x + a 2 x 2 + a3 x 3 + K , in which the a coefficients
are merely the constant terms of the power series. With an initial condition, one can
determine the value of a0 . In this case, y (0) = 1, so y = 1 + a1 x + a 2 x 2 + a3 x 3 + K . After
this step, one determines the polynomial expansion for y , which is simply the derivative
of the polynomial expression already found: y = a1 + 2a 2 x + 3a3 x 2 + L . One can
actually equate these two polynomial expansions in the context of the given differential
dy
equation:
= kx a1 + 2a 2 x + 3a3 x 2 + L = k a 0 + a1 x + a 2 x 2 + a3 x 3 + K
dx
For the sake of simplicity, assume that k = 1. Now one can equate the coefficients on
each side of the equation:
a1 = a0 = 1
1
1
1
2a 2 = a1 a 2 = a1 = (1) =
2
2
2
1
11
1
3a3 = a 2 a3 = a 2 = =
3
3 2 3 2
M
M
M
M
N
1
1
na n = a n 1 a n = a n 1 =
n
n!
2
3
n
x
x
x
Thus, y = 1 + x +
+
+L
= ex.
2! 3!
n!
With these steps in mind, one can solve the following second-order differential
d2y
= ky. This second-order differential equation is important because it is
equation:
dx 2
the kind that describes simple harmonic motion, described earlier in this section. One
must determine polynomial expansions for y, y , and y :
y = a 0 + a1 x + a 2 x 2 + a3 x 3 + L . Given the initial condition that y (0) = 1,
simplified context for the method of power series. Suppose one were to solve

y = 1 + a1 x + a 2 x 2 + a3 x 3 + L.
y = a1 + 2a 2 x + 3a3 x 2 + L.
y = 2a 2 + 6a3 x + L.
These polynomials can be equated to each other in the context of the differential
d2y
= ky (2a 2 + 6a3 x + L) = k (1 + a1 x + a 2 x 2 + a3 x 3 + L).
equation:
2
dx
For the purpose of simplification, assume that k = 1 :
( 2 a 2 + 6 a 3 x + L) = 1 a1 x a 2 x 2 a 3 x 3 L .
Equating the coefficients:
1
2 a 2 = 1 a 2 =
2
a
6a3 = a1 a3 = 1
6
( 1 / 2) 1
12 a 4 = a 2 a 4 =
=
12
24
a3
a1
20 a5 = a3 a5 =
=
20 120

73

 a
1
a
1
y = 1 + a1x x 2 1 x 3 + x 4 + 1 x 5 + L
120
24
6
2
Note that since there is no a1 on the left side of the original equality, yet one on
the right side, a1 must equal zero. Thus,
1
1
y = 1 x2 + x4 +L.
2!
4!
One recognizes this as the power series expansion for y = cos x. In fact, this is what
one would expect from an equation that models simple harmonic motion. Imagine a
pendulum swinging back and forth at small angles. If one were to record its position
versus time, a resulting plot of the data would yield a cosine or sine function.
As always, the method of power series does have certain limitations. This method
can only be used when the power series used are centered at the origin. Also, it can only
be used for what are known as homogeneous differential equations, in which the only
terms are the function in question and a certain number of derivatives, all multiplied by
d2y
constants. For instance, the differential equation
= ky is homogeneous because
dx 2
the only terms are the function in question and its second derivative multiplied by
d2y
= ky + 1 , however, is not homogeneous
constants. The differential equation
dx 2
because it includes a term (1) that is not the function in question or its derivatives
multiplied by constants. A more extensive coverage of this material is the stuff of a
course in ordinary differential equations, but one would ordinarily have to take Calculus
III and Linear Algebra before it.
Concluding Remarks
This chapter introduced the theory and applications of a very important
mathematical tool power series. These expressions are similar to the infinite series
covered in the previous chapter, only power series contain variable terms. In addition,
when determining the convergence of a power series, one must specify where it
converges. Two important types of power series, Maclaurin and Taylor series, were
discussed in the context of polynomial approximations. The technique of polynomial
approximations to difficult functions was also covered in greater depth through a
discussion of their indispensable applications to problems in science and mathematics.
Key Terms:
power series
radius of convergence
interval of convergence
term-by-term differentiation
term-by-term integration
Maclaurin series
Maclaurin polynomial

Taylor series
computational cost
Taylor polynomial
virial equation
Taylors theorem
method of power series
remainder
Lagrange form of the remainder
simple harmonic motion
Lennard-Jones potential

74

Chapter 5: New Coordinate Systems: Parametric

and Polar
This chapter will extend the theory and application of differential and integral
calculus to new coordinate systems. Thus far, all of the mathematics have been done in
the context of what is known as rectangular, or Cartesian, coordinates, in which a
point is assigned an x- and y-value. This coordinate system was devised by the great 17thcentury mathematician and philosopher Ren Descartes. In effect, this novel system
bridged the gap between algebra and geometry, giving rise to the field of mathematics
known as analytic geometry. Unfortunately, Cartesian coordinates are not appropriate
for all problems. It is often quite difficult or impossible to model certain phenomena on
the traditional xy plane. This chapter will introduce two coordinate systems often used
when modeling in Cartesian coordinates fails.
Arc Length
Before continuing on to the study of new coordinate systems, it is important to
introduce an application of the definite integral that, while having no appropriate place in
the previous chapters, is, nonetheless, tested on the AP Calculus BC exam, and is also
applied to the new systems discussed in this chapter This application is the determination
of the length of a curve, or arc length. While the definite integral is useful in
determining the area under the curve on a certain interval, it is equally powerful in
determining the length of the curve on that interval. Geometrically, the determination of
arc length is a one-dimensional application of the definite integral. Finding the volume
of a solid of revolution through the use of the definite integral is a three-dimensional
application because it concerns the summation of an infinite number of three-dimensional
geometric figures (i.e. discs, washers, and cylindrical shells). Finding the area under a
curve is a two-dimensional application of the definite integral because the summation is
of an infinite number of two-dimensional geometric figures (i.e. rectangles, trapezoids,
etc.). Finding the length of a curve is a one-dimensional application of the definite
integral because it concerns the summation of an infinite number of one-dimensional
geometric figures, or line segments. Consider the curve below and the very small
(infinitesimal) line segment associated with it:

The infinitesimal line segment dl is the hypotenuse of the triangle formed by the
perpendicular line segments dy and dx. Using the Pythagorean theorem:
dl 2 = dy 2 + dx 2 dl = dy 2 + dx 2 . Casting this equation in a form that contains the

75

dy 2
dy 2

dx
1 + 2 = 1 + dx . Many

dx
dx

infinitesimal line segments can be summed along a closed interval a x b by using a

definite integral:

dy
recognizable derivative
, dl =
dx

dy 2
1 + dx .
a
dx
The above is the formula for arc length (l). Oftentimes, this integral is far too
complicated to evaluate analytically by conventional methods, so the graphing calculator
is often used for arc length problems. This formula can be used to prove that the
circumference of a circle is indeed 2 r . Consider the formula for a semicircle with a
l=

radius r: y = r 2 x 2 , which is graphed below:

2x
dy
=
=
dx 2 r 2 x 2
One then uses the formula for arc length:

One first determines the derivative:

l=

Use

x
1 +
r 2 x 2

du
a2 u2

= sin 1

dx =

x2
1 + 2
2
r x

dx =

x
r x2
2

r 2 x2
x2
2
+
2
r 2 x2
r x

dx =

r2
2
2
r x

dx

u
+ C. Let u = x, du = dx, a = r :
a
r

r
r
r2

r2
r
x
r
r
2 2 dx =
= rsin1 rsin1 = rsin1 (r) rsin1 (r)
dx =
dx = rsin1
2
2
2
2
r r
r r x
r r x
r r x
r
r
= r . This is the length of the semicircle, so the length of the whole circle must
be 2 r. It appears as though the ancient mathematicians were right.
r

Ex.) Determine the length of the curve y = ln x sin x from x = 1 to x = 2.

First, one should determine the derivative and square this result:
2

ln x cos x sin x sin 2 x

sin x
dy
dy
+
= ln x cos x +
= ln 2 x cos 2 x + 2
.
dx
x
x
x2
dx

76

From this result, one can determine the arc length:

2
dy 2
ln x cos x sin x sin 2 x
2
2
1
=
1
+
ln
cos
+
2
+
dx.
+
dx
x
x

x
x2
a
1
dx
This is a thoroughly unpleasant integral that should be solved on the TI-84.
However, it difficult even to type this expression correctly on the calculator. There is a
trick to overcome this problem. On the Y= screen, type in 1 + ln 2 x cos 2 x for Y1.
After this, scroll down to Y2 and press the VARS key, go to the Y-VARS menu, and
select Function. This should display a list of functions (i.e. Y1 , Y2 , Y3 , etc.) . The first

l=

of these, Y1 , represents 1 + ln 2 x cos 2 x . Select Y1, which should now appear in the Y2
ln x cos x sin x
." After scrolling
slot on the Y= screen. In the Y2 slot, type Y1 + 2
x
down to the Y3 slot, go to the VARS menu as before and select Y2, which should now be
sin 2 x
. Finally,
displayed on the Y3 slot on the Y= screen. In this slot, type Y2 +
x2
scroll down to the Y4 slot and press the 2ND and x2 buttons, which should put a radical
symbol in the Y4 slot. Go to the VARS menu, selecting Y3. This should put the Y3
within the radical symbol in the Y4 slot. The function Y4 actually represents
ln x cos x sin x sin 2 x
1 + ln x cos x + 2
. Overall, the Y= screen should look like
+
x
x2
this:
2

Use the method explained above whenever a structurally difficult problem must
be solved graphically.
Now, one can evaluate the integral of Y4 from x = 1 to x = 2 :
2

ln x cos x sin x sin 2 x

l=
+
dx = 1.1965149.
1 + ln x cos x + 2
x
x2
1
Arc Length With Vertical Tangents: Sometimes, the derivative of a function for which
one is determining the arc length is not integrable at a certain point on the interval in
question. There are ways of remedying this, however. In the case of a vertical tangent,
one must find the arc length of the inverse of the original function, meaning that the
equation and the limits of integration must be in terms of y instead of x. The inverse will
have the same length as the original function, but the positions of the x and y will be
reversed.

77

Ex.) Find the length of the curve y = 5 x on the interval [-2,3].

1
dy
. This derivative has a
=
5
dx 5 x 4
vertical tangent at x = 0. Thus, one takes the inverse of the original function and finds its
derivative:
dx
x( y ) = y 5
= 5y4.
dy
Taking the derivative of the function,

l=

dx 2
1 + dy ( Where c and d are the y - values associated with a and b).
dy

l=

1 + 25 y 8 dy. Note that in the calculator, this can be put in as 1 + 25 x 8 ,

since the calculator does not know the difference between y(x) and x(y). Solving with the
graphing calculator, one gets l = 6.3083997.
Arc Length With Corners or Cusps: Recall that a function is not integrable if there are
corners ( ) or cusps ( { ) along the interval of integration. Thus, if the derivative of
the function whose arc length one is trying to find has any of these points, one must
work around them by making a piecewise function. Consider, for instance, the
5
2
5 3
3
function y = x . The derivative of this function is x , which has has a cusp at x = 0 :
3

Thus, the original function must be defined as a piecewise function about the
53

origin: y = x 5 , x 0 One must then integrate the first piece from the lower limit to 0
x 3 , x 0.
and then integrate the second piece from 0 to the upper limit.
5

Ex.) Determine the length of the curve y = x 3 on the interval [-2, 2].
One integrates each piece separately and then adds each numerical result
to yield the total arc length:

78

l=

5 23
1 + x dx +
3

5 23
1 + x dx = 3.8493405 + 3.8493405 = 7.698681.
3

Introduction to Parametric Curves

If one were to throw a ball in the air and track its horizontal and vertical position
with respect to time, it would be difficult to analyze this situation in a Cartesian
coordinate system. Instead, one wishes to define both the horizontal position x and the
vertical position y in terms of time t. One would then have two functions defined by the
same parameter t: x = f (t ) and y = g (t ) . Thus, at any time t, one can determine both
the horizontal and vertical positions of the ball. The two equations that are defined in
terms of the same parameter are known as parametric equations and, when graphed,
yield parametric curves. Parametric curves are very helpful in determining the path
associated with an objects motion. While functions such as x(t) and y(t) are helpful in
determining aspects of motion such as position, velocity, and acceleration, as studied in
AP Calculus AB, they cannot directly convey what that motion looks like. Consider,
again, the example of throwing a ball in the air. In this situation, neglecting air
resistance, the position functions are as follows:
x(t ) = a + bt and y = c + dt et 2 , where a, b, c, and d are all constants.
However, these equations convey nothing directly about the path of the particle,
which is y and function of x (i.e. y(x).). The parametric equations allow one to define a
coordinate point (x, y) in terms of the parameter: (x(t ), y (t ) ) .
On the graphing calculator, one can display parametric curves by going to MODE
and down to PAR rather than FUNC. In the Y= screen, each region corresponds to two
parametric equations, x(t ) and y (t ). In the X1T slot, put 1 + t and in the Y1T slot, put
1 + t t 2 (Note that the calculator uses a capital t to indicate a parameter). The resulting
graph should look like this:

This is exactly the path that one would expect from throwing a ball in the air a
parabola. But why is nearly half of the parabola missing in the graph? This concerns a
concept known as the parameter interval, which is the domain of the parameter. The
calculator has, in the case of the curve above, automatically made this interval
0 t 2 . Thus, the initial point of the parabola occurs at (x(0), y (0) ) and the
terminal point of the parabola occurs at (x(2 ), y (2 ) ). In this case, one does wish for
the initial point to occur when t = 0 because t is the time parameter. In some cases,
however, it may be necessary to define a specific parameter interval. One can change the

79

parameter interval on the graphing calculator by going to the WINDOW menu while in
parametric mode.
Converting Between Parametric and Rectangular Form: Sometimes, it is
desirable to convert from parametric form back to rectangular (Cartesian) form. While
this can often be done algebraically, there are many cases in which this is very difficult.
The parametric equations already discussed are an easy example. Given x(t ) = 1 + t and
y (t ) = 1 + t t 2 , one can find y as a function of x:
x = 1+ t t = x 1
y = 1 + t t 2 y = 1 + ( x 1) ( x 1) 2 = 1 + ( x 1) ( x 2 2 x + 1) = x 2 + 3 x 1.
It is almost certainly the case that one will not be expected to derive parametric
form from rectangular form, as this often requires either a lot of mathematical theory, or
a lot of mathematical experiment. There are some certain cases in which this is not that
difficult, but it is not expected on the AP Calculus BC exam. On the exam, one will be
given parametric equations and asked to perform an operation on them or apply them in a
certain way. Not only will one not be asked to convert from rectangular form to
parametric form, but it is also very rare that one will be asked to convert from parametric
form to rectangular form. Below are some examples of common curves seen in
parametric form:
The Circle: This is the only curve for which a derivation of the conversion
between Cartesian and parametric coordinates will be shown. Consider a particle moving
along the circle below:

The parameter t represents an angle in this case. One can define the particles xand y-positions in terms of this parameter through the use of right triangle trigonometry.
If a is the hypotenuse:
x = a cos t and y = a sin t.
Thus a circle has the parameterized form of x(t ) = a cos t , y (t ) = a sin t , where a is
the radius of the circle. For instance, the parametric equations x(t ) = cos t , y (t ) = sin t
represent the unit circle.
The Ellipse: An ellipse has the parameterized form x = a cos t , y (t ) = b sin t. If
a > b, then a is referred to as the major axis and b is referred to as the minor axis. If
a < b, the reverse is true. For instance, the ellipse x(t ) = 2 cos t , y (t ) = sin t is shown
below:
80

The major axis is 2 and the minor axis is 1. Note that a circle is merely an ellipse
with a = b.
The Astroid: At this point, more esoteric parametric curves will be considered, but
they all may very well appear on the AP Calculus BC exam. An astroid (not asteroid)
is a parametric curve generated by tracing the path of a point at the edge of a circle
rolling inside of another circle. The general parametric form of an astroid is
x(t ) = a cos 3 t , y (t ) = a sin 3 t , in which a is the radius. For instance, the astroid
x(t ) = cos 3 t , y (t ) = sin 3 t is shown below:

The Cycloid: A cycloid is a parametric curve defined by the path of the point on
the edge of a circle that rolls along a straight line. It has the general parametric form
x(t ) = a(t sin t ), y (t ) = a (1 cos t ). Note that in the cycloid, unlike in other situations,
the x-coordinate is associated with the sine function and the y-coordinate with the cosine
function. The cycloid played a prominent role in physics during the time of Newton
because it was used to find solutions to the brachistochrone and tautochrone problems.
The former concerned the curved path between two points in which an object covered the
distance of the curve in the least amount of time, and the latter concerned the curious
observation that an object in curved space, such as a bowl, will take the same amount of
time to reach the bottom regardless of its starting point on the curve. As an example,
consider the cycloid x(t ) = 3(t sin t ), y (t ) = 3(1 cos t ) shown below:

81

The Witch of Agnesi: The Witch of Agnesi is an example both of parametric form and of
linguistic error. This curve was discovered by the 18th-century Italian mathematician
Maria Gaetana Agnesi. It is described by the parametric equation
x(t ) = 2 cot t , y (t ) = 2 sin 2 t , which shown graphically below:

The name witch of Agnesi is a misnomer. In Italian, the phrase is la versiera di

Agnesi, in which la versiera refers to a curve. In 1801, however, Cambridge
University mathematician John Colson mistranslated this part of Agnesis work as
lavversiera, which means wife of the devil. Like many other linguistic foul-ups in
history, this translation has stuck.
The Calculus of Parametric Curves
The algebra and geometric interpretation of the derivative and the integral are
different for parametric curves in comparison to the Cartesian curves studied hitherto this
chapter. In this section, differentiation and integration will be performed upon parametric
equations and the applications of these operations will be discussed in terms of what has
already been learned about derivatives and integrals.
Differentiation of Parametric Equations: The derivative in parametric language still has
dy
. However, recall
the same form as in Cartesian language; it has Leibnizs form of
dx
dy
that y and x are now defined in terms of a parameter t. In effect,
actually represents
dx
the quotient of the derivative of y with respect to t and the derivative of x with respect to
dy
dy dt
t:
. This formula is derived through the use of the chain rule for derivatives:
=
dx dx
dt
dy dy dt
=
. Similar to the derivative in Cartesian coordinates, the derivative in
dx dt dx
parametric coordinates can be interpreted as the slope of the tangent line to a curve.
Ex.) Determine the slope of the tangent line to the ellipse defined by
2
x = 2 cos t , y = 5 sin t when t =
.
3

82

dy
2
dy
dy dt 2 sin t
=
=
= tan t
5 cos t
5
dx
dx dx
dt

t=

2
3

2 2
= tan
= 0.692820323.
5 3

The slope of the tangent line to a parametric curve has an important interpretation
in physics. In the context of particle motion, the derivative of a parametric equation
represents the direction in which the particle is moving. For instance, a derivative of zero
implies that the particle is moving horizontally and a derivative of infinity implies that
the particle is moving directly up or down.
Ex.) A particle is moving in two dimensions as described by the parametric
equations x(t ) = sin t , y (t ) = 2t 2 5. At t = 3.5, what is the slope of the tangent line and
in which direction is the particle moving?
dy
dy
dy dt
=
= 4t sec t
= 4(3.5) sec(3.5) = 14.94997066. The particle
dx t =3.5
dx dx
dt
is moving upwards to the left.
To see why this is the case observe the parametric curve below, which represents
the particles path:

The black dot represents the particle, which is moving upwards to the left.
The evaluation of the second derivative of a parametric equation is slightly more
difficult than finding the first derivative. The second derivative ( d 2 y / dx 2 ) is not merely
the quotient of the second derivatives of the x- and y- equations with respect to t. Instead,
one must apply the chain rule in a somewhat less obvious way:
dy
d
dx
2
d y d dy d dy dt
= =
= dt . Here is an algebraically difficult
2
dx
dx dx dt dx dx
dx
dt
problem:

83

Ex.) Determine the second derivative of the Witch of Agnesi.

First, determine

dy
:
dx

dy
2 sin t cos t
dy dt 4 sin t cos t
=
=
=
= 2 sin 3 t cos t.
2
1
dx dx
2 csc t
dt
sin 2 t

Next, determine

d dy
:
dt dx

d dy
3
3
= 2 (sin t )( sin t ) + (cos t )(6 sin t cos t )
dt dx
Finally, divide this expression by

dx
:
dt

dy
d
dx
3
3
3
3
2
dt = 2 (sin t )( sin t ) + (cos t )(6 sin t cos t ) = (sin t )( sin t ) + (cos t )(6 sin t cos t ) = d y .
dx
dx 2
csc 2 t
2 csc 2 t
dt

] [

Higher-order parametric derivatives are not tested on the AP Calculus BC exam.

However, they follow the same general pattern as the second derivative. For instance, the
d n 1 y
d2y
d 2
d n 1
dx
dx
3
n
d y
d y
dt
dt
=
=
. In general,
.
third derivative has the form
3
dx
dx
dx
dx n
dt
dt
Integration of Parametric Equations: For a curve y(x) parameterized by t, integration is
as simple as integrating both x and y with respect to t. The methods for evaluating these
integrals are the same as those used for evaluating integrals in Cartesian coordinates. On
the AP Calculus BC exam, one is not usually asked to simply evaluate the integral of
parametric equations, but to apply this integration to a certain problem. On the exam,
only one application is tested arc length, which was discussed earlier in this chapter in
the context of Cartesian coordinates.
Arc Length of Parametric Curves: Recall that the formula for arc length is
dy 2
1 + dx . For parametric equations, this formula must be written in terms
a
dx
of the parameter t. This is accomplished by using the formula for the first parametric
derivative:

l=

84

l=

dy 2
1 + dx =
dx

b
a

t2
t1

b
a

dy / dt 2
1 +
dx =
dx / dt

(dx / dt ) 2 + (dy / dt ) 2 dx
dt =
dt
(dx / dt )
2

t2

t2
t1

b
a

(dx / dt ) 2 + (dy / dt ) 2 dx
dt

(dx / dt ) 2

dt
2

dx
dy
+ dt
dt
dt

dx
dy
l =
+ dt , where t1 is the value of t associated with a and t 2 is the
dt
dt
t1
value of t associated with b. Note that in the context of particle motion, this
parameterized formula refers to the total distance traveled by the particle.

Ex.) The position function r (x(t ), y (t ) ) of a particle is an ellipse given as

x(t ) = 6 cos t , y (t ) = 3 sin t. Determine the total distance traveled by the particle from
t = 0 to t = 4.
2

4
dx
dy
( 6 sin t )2 + (3 cos t )2 dt.
+
dt
=


dt
dt
t1
0
Evaluating with the graphing calculator, which should be done in FUNCTION

l=

t2

mode,
l=

36 sin 2 t + 9 cos 2 t dt = 17.783441. Note that the integral could have been

easily evaluated analytically, but one would had to have used the calculator eventually
when plugging in a value of 4 for t.
Area Under Parametric Curves: It is hardly ever the case that one is asked to
determine the area enclosed by a parametric curve, and certainly not on the AP Calculus
BC exam. However, since this information is so scant, I thought it appropriate to include
it in this book. The derivation is really quite simple. Recall that the area under a curve
x2

y( x) dx. When this curve is

dx
parameterized by t, the integral must be evaluated with respect to t:
y (t ) dt. This
dt

y(x) from x1 to x 2 is equal to the definite integral

x1

t2

t1

last formula represents the area under a parametric curve from t1 to t 2 .

Ex.) Determine the area enclosed by one arch of the cycloid represented by
x = 3(t sin t ), y (t ) = 3(1 cos t ).
One can determine the limits of integration by observing which two first
values of t correspond to y-values of zero. This occurs at t = 0 and t = 2 .

85

Area =

(3 3 cos t )(3 3 cos t )dt =

(9 cos 2 t 18 cos t + 9)dt.

1 + cos 2
. Remember that the common
2
trigonometric identities must be memorized for the AP Calculus BC exam. They can be
found in appendix A of this book.
2
2
2
27
9 9 cos 2t

9
18 cos t + 9 dt =
(9 cos 2 t 18 cos t + 9)dt =
+
cos 2t 18 cos t + dt
2
2

0
0 2
0 2
Use the trigonometric identity cos 2 =

27
1
= sin t 18 sin t t = 27 .
2
2 0

Surface Areas of Revolution of Parametric Curves: A very common application of

integration of parametric equations, and yet one that is not covered in the AP Calculus
BC curriculum, is the determination of the surface area of a solid generated by revolving
a bounded parametric curve about an axis. Note that this is the surface area or
revolution, and not the volume, as is usually the case in applications of the definite
integral. While the derivation will not be shown here, the formulae for the surface areas
of revolution for parametric curves are given below:
(When revolved about the y-axis): Surface Area =

dx
dy
2 y + dt
dt
dt
a

dx
dy
2 x + dt .
dt
dt
a
Notice that the radical expression is the arc length, which is multiplied by another
linear term to yield the surface area.
(When revolved about the x-axis): Surface Area =

Here is a rather strange problem concerning surface areas of revolution of parametric

curves:
Ex.) A baker, striving to be as efficient as he possibly can, wishes to determine the
total area of icing needed to completely cover each of 25 doughnuts with theoretically no
icing left over. A mathematically savvy baker, he decides to develop a mathematical
model of his doughnuts which consists of a circle defined parametrically as
x(t ) = 2.3 cos t + 4, y (t ) = 2.3 sin t + 4. This circle is revolved about the x-axis to generate
the model of a given doughnut. What is the approximate area of total icing needed to
completely cover 25 doughnuts with no icing left over?
Since the axis of revolution is the x-axis, the formula used would be
b

dx
dy
2x + dt . What are the limits of integration? Starting at t = 0, one
dt
dt
a
reaches this point again at t = 2 . Thus, these are the two limits of integration.

86

Surface Area =

2
0

2 (2.3 cos t + 4) ( 2.3 sin t )2 + (2.3 cos t )2 dt

Using the graphing calculator in FUNCTION mode, one can easily evaluate this integral:
Surface Area = 157.91367. This is the surface area of one doughnut.
Multiplying by 25 yields the total area of icing that the baker needs:
Total Surface Area = 25157.91367 = 3947.84175 square units of icing.
Introduction to Polar Curves
Unlike Cartesian and parametric coordinates systems, there is a system
that defines points in space not by x- and y- coordinates, but by lines and angles. This
coordinate system is associated with polar equations. A point in space in a polar
coordinate system is defined by a directed distance (r) of certain length and the angle
(in radians) at which this distance is displaced from a reference ray, known as the initial
ray. Thus a point P has the coordinative notation of P(r , ) . This is supported by the
diagram below:

Whereas in Cartesian coordinates the point of reference is referred to as

the origin, in polar coordinates, the point of reference is referred to as the pole. Also,
unlike in Cartesian and parametric systems, points in polar coordinates are not unique.
That is, while a certain r and combination may define a certain point in polar space,
there are other values of r and that define that same point. For instance, consider the
point defined as (3, / 4). There are other (actually infinitely more) combinations that
define this same point, like (3, 9 / 4) and (3, 5 / 4). While the first one seems obvious
( 9 / 4 is 360 away from / 4) the second one needs more explanation. What is the
relationship between 9 / 4 and 5 / 4 ? They are separated by 180, or radians.
Making the 3 negative is essentially the same as moving it 180. To understand this
better, observe the diagram below:

87

In order to have the directed distance that is 9 / 4 radians away from

the initial ray (i.e. the positive x-axis) it would be necessary to negate the 3, effectively
moving the directed distance 3 units backward.
Note also that radians measured from the counterclockwise direction
are positive and radians measured from the clockwise direction are negative. For
instance, the polar coordinates (3, / 4) define the same point as (3, 7 / 4) . This is
because 45 is equivalent to -315.
Converting Between Polar and Cartesian Coordinates: By observing the first supplied
diagram of this section, it should be obvious that polar coordinates and Cartesian
coordinates are related by right-triangle trigonometry. The x-value is related to r and
by the cosine function and the y-value by the sine function:
x = r cos
y = r sin .
One can also convert from Cartesian coordinates to polar coordinates. To do this, one
simply uses the Pythagorean theorem and right-triangle trigonometry:

r = x2 + y2
y
= tan 1 , x 0.
x
Ex.) Express the equation xy = 4 in polar form.
Since x = r cos and y = r sin , (r cos )(r sin ) = 4
r 2 cos sin = 4.
Polar Curves: Polar curves are defined as the directed distance r as a function of (i.e.
r ( ) . Thus, as r and vary, they define a curve in polar space. For instance, consider
the simplest polar curve, a circle. A circle has a constant r with an angle that varies.
Thus, the polar equation for a circle is simply r ( ) = constant, where the constant is the
radius. To graph in polar form on the TI-84, go to MODE and choose POL. In the Y=
screen, one should now see a list of r =. To graph a circle with a center at the origin
and a radius of 3, simply input 3 for r1. The result is:

88

Nowadays, most polar curves are electronically generated on the

graphing calculator. However, not so long ago, students used polar graph paper to draw
polar curves. Polar graph paper looks quite different from regular graph paper:

It will almost certainly not be required to produce a polar curve on polar graph paper
on the AP Calculus BC exam. Even if one is asked to reproduce a polar curve on paper,
it will most likely be in the context of regular graph paper.
Obviously, polar curves become more elaborate once one allows r to vary along with
. Below is a discussion of the various sorts of polar curves that one is most likely to
encounter:
Spirals: In polar coordinates, a spiral can be viewed as a curve that turns around a
point of reference while getting closer or farther away from it, depending upon the
equation representing the spiral. Ancient mathematicians were fascinated by these
geometric phenomena because they are often manifested in nature. Many of the spirals
discussed here were first used to model these natural systems.
The Archimedean spiral has the general form r ( ) = a + b , where a and b are
constants. For instance, the Archimedean spiral r ( ) = 1 + 2 is shown below:

89

Note that the calculator automatically makes the interval of the angle [0, 2 ] . To
view more of the spiral, one could make max larger in the WINDOW option. For
example, allowing max to equal 4 yields a fuller spiral:

A key feature of the Archimedean spiral is that the distance between successive turns
is always the same. This is not the case for the logarithmic spiral, which has the general
form r ( ) = ae b , where a and b are constants. For instance, consider the logarithmic
spiral r ( ) = e

For this graph max was made 8 . The turns near the origin are so small compared to
the noticeable turns that they only appear as a dark spot. The primary observation to take
away from this is that at each subsequent turn, the distance of that turn relative to the one
before it increases exponentially. Logarithmic spirals are extremely prevalent in nature.
They are seen, for instance, as mollusc shells, spiral galaxies, hurricanes, and the cochlea
(an organ of the inner ear) of mammals. There are many other spirals that can be
expressed in polar form, but they are slightly more esoteric and will not be covered here.
The Limaons: A limaon is a polar curve with the general form r ( ) = a b sin or
r ( ) = a b cos , where a and b are constants. The term comes from the Latin limex,
which means snail, referring to the general shape of these polar curves. The limaon
was first described by the German Renaissance artist Albrecht Drer in his
Underweysung der Messung (Instructions in Measurement) (1525). The limaon family
is generally composed of four types of curves, which are shown below:

90

convex limaon
Conditions: 2a b

limaon with inner loop dimpled limaon

Conditions: a < b
Conditions: a < b < 2a

cardioid
Conditions: a = b

Note that memorization of these terms is not required for the AP Calculus BC exam,
but it is helpful to know them nonetheless. The cardioid is so well-studied among the
limaons that it has its own name! It is a very common polar curve in which a = b.
Geometrically, a cardioid looks similar to a heart shape, hence its name. The cardioid is
generated by tracing the point on the edge of a circle as it rolls along a circle of the same
radius. For instance, consider the cardioid represented by r ( ) = 2 + 1.5 sin :

In microphone science, the cardioid is the most common shape for the sounds picked up
by a unidirectional microphone. A unidirectional microphone is sensitive only to sounds
coming from a specific direction, which is helpful when one does not want to pick up
extraneous noise.
Besides the cardioid, the AP Calculus BC exam often includes the limaon with a
loop, in which a < b . For instance, the limaon represented by r ( ) = 1 + 2 cos is
shown below:

As will be discussed in upcoming section, the AP exam often likes to have students
analyze the loop in some way.
The Lemniscate: A lemniscate is geometrically a figure-eight. It was first described by
Jakob Bernoulli (a member of that famous mathematical dynasty) in 1694. The name is
derived from the Latin lemniscus, which means ribbon. It is also sometimes used to

91

describe the symbol for infinity ( ) . It is described by the general polar equation
r 2 = a 2 cos 2 , where a is a constant. Unfortunately, the lemniscates cannot be graphed
in this form on the TI-84; one must put the equation in as r ( ) = a cos 2 . Thus, in
one slot on the Y= screen, one would input the positive value, and in another slot the
negative value. For instance, the lemniscate represented by r 2 = 25 cos 2 would be
written as r ( ) = 5 cos 2 , which is shown below:

Unfortunately, this is the best the calculator can do. The space where the calculator
thinks no curve exists is due to taking the square root of a negative number at those values, yielding imaginary numbers. Nevertheless, the curve really does exist in this
space; it is just that the equation had to be put into a form that the graphing calculator
would recognize.
The Rose Curves: The final class of polar curve discussed in this section is the rose
curve. As suggested by its name, a rose curve consists of petals that emanate from a
single point. Rose curves has the general form r ( ) = a cos(b ) or r ( ) = a sin(b ),
where a and b are constants. The constant b is rather significant because it determines
how many petals the rose curve will have. If b is an even number, the number of petals
is 2b, while if b is an odd number, the number of petals is simply b. For instance,
consider the rose curve represented by r = 2 cos(6 ) :

Since b is an even number (6) the number of petals is twice b (12). What if b is an
irrational number like ? The number of petals will be irrational. However, if one
allows the limits on the magnitude of to be great enough, one would see a disc
generated. Consider the rose curve represented by r ( ) = cos( ) :

92

This is not a complete curve because the number of petals is irrational. The arrow
indicates where the curve stops. If one goes to WINDOW and makes max very large
(perhaps 20 ), the resulting graph looks a lot different:

If one thinks about this in terms of the philosophy of calculus, it is theoretically the
case that an irrational number as b results in an infinite number of petals. If one could
truly make = , the graph would look completely solid; it would be a complete disc.
Determining Points of Intersection: Solving polar equations is very important for the
next section of this chapter. It is not terribly difficult, but it requires more thought than in
solving equations in Cartesian or parametric coordinates. In general, one can decide
where two polar curves intersect by setting their polar equations equal to each other and
solving for . However, recall the property of polar coordinates that a -value
associated with a point is not unique; there are infinitely many other -values that are
associated with this same point. Thus, just because the -value of one polar curve at a
certain value of r is not the same as the -value of another polar curve at the same value
of r does not mean that the curves do not intersect. Setting the equations of the two polar
curves equal to each other yields the points where the curves intersect due to having the
same r- and -value. However, consider the two curves r1 ( ) = sin and
r2 ( ) = 1 sin . From the diagram below, there seems to be three points of intersection:

93

Setting the two equations equal to each other, one can find the -values at which the
curves intersect:
1
r1 ( ) = r2 ( ) sin = 1 sin 2 sin = 1 sin =
2
1 5
= , .
6 6
Thus, the two coordinates (r , ) where the two curves intersect are:

1
1 5
, and , .
2 6
2 6
But did it not seem, by looking at the graph, that there were three points of
intersection? If one were to zoom in on the graph on the TI-84, it would be clear that
both curves occur at the origin and, thus, intersect. Thus, there are three points of
1 1 5
intersection: , , , , and (0,0). Unfortunately, solving a trigonometric equation
2 6 2 6
will not confirm this because the solutions to these equations only show the points at
which both components of the coordinates are equal, and, as it turns out, while the rvalue of the coordinate of each curve at the origin is the same (0), the -value differs.
Solving for when r = 0 in each equation,
0 = sin = 0
0 = 1 sin =

Thus, at the origin, the coordinates for r1 are (0, 0) , while for r2 they are 0, .
2
However, these both represent the same point, the origin! The main message to take
away from this is that solving the resulting trigonometric equation from setting two polar
curves equal to each other might not convey all of the points of intersection because there
could very well be points of intersection in which the coordinates of the two curves are
not the same. The best way to check if this is the case is to analyze the actual graph on
the calculator.
The Calculus of Polar Curves
In this section, one will see a number of similarities between polar calculus and
parametric calculus. While the methods of differentiation and integration may seem
different from those encountered before, they still have the same interpretations and
applications; the derivative still refers to the slope of the tangent line and the integral is
still used to find area, arc length, or properties of solids of revolution.
Differentiation of Polar Equations: The formula for the derivative of a polar equation is
derived in the same manner as that applied to parametric equations; the chain rule is used.

94

dy

The derivative still has Leibnizs notation i.e. , but equations in polar form have x
dx

and y (together as r) in terms of . Thus, one applies the chain rule for derivatives:
dy
d
(r sin )
dy dy d d
d

=
=
=
.
dx
d
dx d dx
(r cos )
d
d
Applying the product rule to the numerator and denominator of the final expression
above:
d
dr
(r sin ) r cos +
sin
dr
d
d

is easy to evaluate
. Note that the expression
=
dr
d
d
(r cos )
cos r sin
d
d
because it represents the derivative of the polar curve with respect to , which does not
dr
does not represent the
require any algebraic manipulation. Be wary, however, that
d
dy
represents this. The AP Calculus BC
slope of the tangent line to the polar curve; only
dx
exam often asks questions pertaining to vertical and horizontal tangents in regards to the
dy
dy d
differentiation of polar curves. With the formula
in mind, a polar curve will
=
dx dx
d
dx
dy
dy
= 0 and
0 and a horizontal tangent when
=0
have a vertical tangent when
d
d
d
dx
and
0. Note that if both derivatives are zero in any of these cases, an indeterminate
d
0
form of results and no conclusion can be drawn.
0

Ex.) Find the horizontal and vertical tangents of to the limaon represented by
r ( ) = 3 4 cos on the interval [0, 2 ].
To find any vertical tangents, one sets

dx
dy
to zero, also making certain that
does
d
d

not equal zero:

dx
d
dr
r cos =
=
cos r sin = 4 sin cos 3 sin + 4 cos sin = 8 cos sin 3 sin = 0
d d
d
It would be a good idea to solve this equation via calculator by going back to
FUNCTION mode and treating the s as if they were xs . In this way, one can
graphically determine where the expression equals zero. Doing this, one should get:
95

 = 0,1.863996, , 5.0967858, 2 .
One must now make certain that these values do not yield a value of zero for

dy
:
d

dy
d
dr
r sin = r cos +
=
sin = 4 sin 2 3 cos 2 3 cos . Testing all of the
d d
d
dy
values of found above, one would find that none of them make
zero. Thus, one
d
can be certain now that there are vertical tangents at = 0,1.863996, , 5.0967858, 2 .
To find any horizontal tangents, one would set

dy
equal to zero and make certain that
d

dx
does not also equal zero:
d
dy
= 4 sin 2 3 cos 2 3 cos = 0. Again, one should solve this equation via
d
calculator by graphing in FUNCTION mode and finding where the graph crosses the xaxis on the interval. Doing this, one should get: = 1.6771037, 4.6060816.
dx
. Testing the two
One must now make certain that these values do not also make
d
dx
dx
, it is found that none of them make
zero. Thus, one can now be certain
values in
d
d
that there are horizontal tangents at = 1.6771037, 4.6060816.
One could check these answers by analyzing each of the discovered -values on the
actual graph of the polar equation.
Integration of Polar Equations: This section will cover the same three applications of
polar equations as was discussed in the section on parametric equations area, arc length,
and surface area of a solid of revolution. However, only the first two are tested on the AP
Calculus BC exam in regards to polar equations.
Area Under and Between Polar Curves: Due to the special characteristics of polar
coordinates in comparison to Cartesian coordinates, the formulae for area under and
between polar curves must be derived geometrically. Consider a small section of a polar
curve as shown below:

96

The diagram above represents a polar curve r = f ( ) that has an area bounded by the
angles and . Within this area is a circular sector with length r (or f ( ) ) , which
sweeps out the angle . Recall from circle geometry that the area of a circular sector
1
has the following formula: A = r 2 , where r is the radius and is the angle swept out
2
in radians. In accordance with the central theme of calculus (see the prelude), one moves
to the level of the infinitesimal and allows the number of sectors within this area to
become infinitely small, such that when these very small sector are added together, the
result is the value of the bounded area:
i =n

A = lim

i =1

1 2
1
r i =
2
2

r 2 d .

This integral does not look very complicated, and, indeed, it is not. However, perhaps
the most complicated part of solving problems concerning the area under a polar curve is
the determination of the limits of integration. Unlike the very simple determination of
the limits of integration in Cartesian coordinates, one must often closely analyze the
graph or solve trigonometric equations when dealing with polar curves. The following
several examples should clarify these necessities.
Ex.) Determine the area enclosed in the cardioid represented by r ( ) = 5(1 + sin ) .
The first step in any problem dealing with area under a polar curve should be an
analysis of the graph:

At the origin, one starts at = 0 . The origin is reached again once the directed
distance r has gone through one full revolution, or 2 . Thus, the limits of integration are
0 and 2 . One can now use the formula for area:
1 2
1 2
2
A=
(5 + 5 sin ) d =
sin 2 + 10 sin + 25 d . One can use the
2 0
2 0

trigonometric identity sin 2 =

1
A=
2

1 cos 2
:
2

1
1 cos 2

+ 10 sin + 25 d =

2
2

1
51 1
51

10 sin cos 2 + = 10 cos sin +

2
2 2
2 0

97

51
80.11061267. Remember, if using the calculator, always find the area by going
2
to FUNCTION mode.
=

Ex.) Determine the area enclosed by the inner loop of the limaon represented by
r ( ) = e cos .
This problem is slightly more difficult. It requires a close analysis of the way in
which the curve is graphed. One must ask, At which -value does the loop begin? and
At which -value does the loop end? These questions can be answered by observing
the polar curve while it is being graphed. To ease this observation, one can slow down
the rate at which the calculator graphs the curve by decreasing the step size in the
WINDOW option. Zooming in on the part of the curve containing the loop,

The line indicates where the loop begins and the arrow indicates the direction in
which the curve is graphed. The loop starts at = 0 and that same point is reached when
= 2 . Unfortunately, 2 cannot be used as the upper limit of integration because that
would imply that one is taking the area of the entire limaon and not just the loop! To
find the area, one must exploit the property of symmetry in the above graph. Instead of
finding the area of the entire loop, one can find the area of one-half of the loop and then
multiply this value by 2. This is the only way to account for only the loop and not the
rest of the graph. Watch the calculator graph the curve again. The calculator finishes
graphing the first half (the bottom half) of the loop when it reaches the origin. To find
the -value associated with the origin, one plugs in 0 for r in the original equation and
solves for the appropriate value of :
e
(0) = e cos = cos 1 = 0.5251355796.

Thus, the bottom half of the limaon is completed when = 0.5251355796. One can
now find the area of the while loop:
Atotal = 2 A1 / 2loop =

0.5251355796

(e cos )2 d = 0.04985493.

Ex.) Determine the area of one petal of the rose curve represented by r ( ) = 3 sin 5 .
For this problem one must, again, consider the pattern of graphing and the symmetry
exhibited by the curve. The rose curve is shown below:

98

The calculator, of course, begins at = 0. The arrow indicates the direction is which
the calculator graphs the curve. Thus, the calculator finishes graphing the first petal
when it reaches the origin the second time. The -value at which this occurs can be
found by plugging 0 in for r in the original equation and solving for the appropriate value
of :
(0) = 3 sin 5 . This equation will equal zero when 5 = . Thus, =

. Now that
5
the two limits of integration have been found, one can find the area enclosed in one petal
of the rose curve:
1
A=
2

1
(3 sin 5 ) 2 d =
2

9 1 cos10 5 9 9 81
9 sin 2 5 d =
=
12.72345025.
=
22
2 0 2 10
20

It is often necessary not only to find the area under one polar curve, but the area
enclosed by two polar curves. In order to solve the integral that represents this area, one
must find at which -values the two curves intersect, for these will be the limits of
integration. The integral is simply the simply the difference between the definite integral
of the outer curve and the definite integral of the inner curve. For instance, in the
diagram below,

The shaded area can be represented generally as:

1
1
1 2
A=
router 2 d
rinner 2 d =
r2 r12 d .
2
2
2

Ex.) Determine the area outside the circle represented by r ( ) = 2 and inside the
lemniscate represented by r 2 = 3 cos .
One must observe the graph to see which area must be found:
99

The shaded area represents outside the circle and inside the lemniscate. To
determine the limits of integration, one must find the -values at which the two curves
intersect:
2
r2 = r1 3 cos = 2 3 cos = 2 cos = = 0.8410686706. Note that
3
at this value of , only 1/4 of the shaded area has been completed. Thus, the difference
between the definite integrals of the two curves must be multiplied by 4 at the end:
A1 / 4 =

1
2

0.8410686706

(2 3 cos )d = 0.5539306. Note that even though the integral is

negative, areas are always positive. Thus, A1 / 4 = 0.5539306.

A = 4 A1 / 4 = 4(0.5539306) = 2.215722545.
Arc Length of Polar Curves: The derivation for the arc length of a polar curve is the
same as the derivation that was discussed for Cartesian and parametric systems. The
formula is still essentially l =

dy 2
1 + dx . However, polar equations are
dx

dy
in terms
dx
of r and . The reader will recall from the subsection dealing with the differentiation of
dr
sin
r cos +
dy
d
polar equations that the result of this derivation was
. Substituting
=
dx dr
cos r sin
d
this expression into the formula for arc length,
expressed in terms of r and . Thus, one must use the chain rule to express

l=

2

dr

+

cos
sin
r


d
dx . For the sake of space, the algebra leading to a
1 +
dr cos r sin

d

more succinct expression will be omitted, though the reader is encouraged to do it as an

exercise. The more attractive form of the expression is:

100

l=

dr
r +
d .
d
2

Ex.) Determine the length of the first two turns of the logarithmic spiral represented by
3

r ( ) = e .
While the limits of integration could be found via calculator, trigonometric intuition
should convey that two turns means two revolutions, which means 4 . Thus, the
formula for the arc length of this spiral is l =

9
+ e 2 d = 20651.08.
16

Surface Area of Revolution of Polar Curves: One can derive the formulae for the surface
area of revolution of polar curves directly from the formulae for parametric curves:

(When revolved about the y-axis): Surface Area =

dx dy
2 y + dt
dt dt
a

dx
dy
2 x + dt .
dt
dt
a
Recall that the radical expression is the arc length and that x = r cos and y = r sin .
Thus, the equivalent formulae for polar curves are:

(When revolved about the x-axis): Surface Area =

(When revolved about the y-axis): Surface Area =

dr
2 r sin r 2 +
d
d

dr
(When revolved about the x-axis): Surface Area =
2 cos r +
d .
d

Again, note that surface areas of revolution are not covered on the AP Calculus BC exam.

Ex.) Determine the surface area of the resulting solid when the entire rose curve
represented by r ( ) = 6 cos 2 is revolved about the x-axis.
This problem, like many problems involving integrals of polar equations, requires
the consideration of symmetry. This particular rose curve has 4 petals. One can, thus,
determine the surface area of revolution by revolving 1 petal about the x-axis and then
multiplying this result by 4. Notice where the calculator begins graphing the curve:

101

It begins at where the dot is indicated and then moves to the left, but does not return
to its original position until it has graphed the whole curve. Thus, an additional
complexity is added; one must determine the surface area of revolution of half of one
petal, multiply this result by 2, and finally multiply this result by 4. One can find the
upper limit of integration by determining the first -value at the origin:
(0) = 6 cos 2 . This expression is zero when 2 =
Surface Area1/2 Petal =

. Therefore, =

(6 cos 2 )2 + ( 12 cos 2 )2 d = 6.7082039.

Surface Area1Petal = 2 Surface Area1/2 Petal = 13.41640787.

Surface Area4 Petals = 4 Surface Area1Petal = 53.66563146.
Concluding Remarks
This chapter discussed a new application of the definite integral, the determination
of arc length, and two new coordinate systems parametric and polar. These coordinate
systems were derived from the Cartesian system through the use of certain laws of
algebra and trigonometry. In addition to the discussion of these systems and the various
equations and curves associated with them, their importance in differential and integral
calculus was also discussed. While the derivative and the integral had the same overall
meaning and applications in parametric and polar curves as in curves in the Cartesian
system, the algebraic derivations of the formulae for these meanings and applications
differed significantly from the Cartesian system. In the next and final chapter on vector
algebra and calculus, the reader will encounter a number of the expressions introduced in
this chapter.
Key Terms:
rectangular (Cartesian) coordinates
analytic geometry
arc length
parameter
parametric equations
parametric curves
parameter interval
initial point
terminal point
astroid
cycloid
Witch of Agnesi
polar equations
directed distance

initial ray
pole
polar curves
Archimedean spiral
logarithmic spiral
limaon
cardioid
lemniscate
rose curve

102

Chapter 6: Vectors and Vector Calculus

The final chapter of this book introduces an immensely important topic in
mathematics and science. Thus far, the problems in this book have dealt with quantities
that are characterized solely by their magnitude (i.e. size). Such quantities are referred to
as scalar quantities. Many quantities in pure and applied mathematics and science,
however, are characterized not only by their magnitude, but also by their direction in
space. Such quantities are referred to as vector quantities, which will be discussed in
this chapter.
Let the reader be aware that this chapter covers much more than is required on the AP
Calculus BC exam. In fact, it is debatable as to whether the AP exam actually tests
vectors at all, even when it claims that it does. When the exam tests on vectors, it is
actually testing on parametric equations and asking the student to use vector notation,
which will be described in this chapter. In truth, vector calculus is technically the first
topic usually covered in a college Calculus III class, for it bridges the gap between singlevariable and multivariable calculus. I feel that I cannot do justice to the rich topic of
vectors without dedicating more energy to it in this book than the AP exam does.
Nevertheless, for students who are using this book for the sole purpose of an AP study
guide, I will indicate the material that will be tested on the exam.
Introduction to Vector Quantities and Their Algebra
When one describes his or her weight, it is not often that you hear him or her say
something like, 140 pounds directed at the negative vertical axis. However, weight is a
force (i.e. the force of gravity), which, like many other quantities in nature, must be
described by its magnitude and direction in space. Force, therefore, must be represented
by a vector. While a vector has a somewhat different appearance in the field of linear
algebra, this is a calculus book, and, thus, the definition used in calculus and physics will
be discussed here. In this context, a vector may represented as an arrow whose length is
proportional to the magnitude of the quantity that it represents and that has a certain
orientation in space relative to some reference frame, usually a set of axes. For instance,
suppose a pilot needs to know the vector that represents the wind velocity relative to the
ground (As will be discussed later in this chapter, position, velocity, and acceleration are
all vector quantities). Suppose that this velocity is 30 km/h, 25 south of north. Also,
suppose that the planes velocity relative to the ground is 8.0 10 2 km/h. One can
represent these vectors pictorially on a set of axes in which the positive y-axis represents
north, the negative y-axis represents south, the positive x-axis represents east, and
the negative x-axis represents west:

103

Unfortunately, these vectors are not drawn exactly to scale; the vector for the
airplanes velocity would be longer than it is in the diagram if 30 km/h is actually
represented by the length it is given. Nevertheless, this diagram should convey the
essential features of vectors, that they depend both upon magnitude and direction.
Whenever describing a vector, one must always take both of these properties into
account. It is also important to note an important property of vectors; they can be moved
anywhere in space as long as their magnitude and direction are preserved. The positions
of the vectors used in the example are not unique. The smaller vector could have been
below the larger vector, the tail of the larger vector could have been drawn at the
head (i.e. the arrowhead) of the smaller vector, or any possible combination, as long
as both magnitude and direction are preserved.
Note that the notation for the vectors given in the previous example was quite
cumbersome. Scientists and mathematicians have developed more succinct ways to alert
a reader to a vector quantity, and sometimes, if one is only working in one dimension
(say, with a positive and negative axis), it usually suffices to describe the direction of the
vector quantity by its sign (i.e. negative or positive). For instance, suppose a particle is
only moving along the x-axis. This particles velocity, for instance, can be indicated by
the size of the number and a positive or negative sign, depending upon where the particle
is moving relative to the origin. In printed texts, vectors are usually indicated by boldface type. In handwritten problems, however, this is impractical, so one usually indicates
v
a vector by placing a small arrow above the symbol that represents a vector (e.g. V in a
handwritten problem would usually be seen as V in a printed text). More vector notation
will be discussed later in the chapter.
Vector Addition: Common operations upon vector quantities such as addition,
subtraction, multiplication, or division are carried out very differently from scalar
quantities. The primary reason for this is, of course, that one must take into account the
vector quantitys direction as well as its size. Addition and subtraction will be discussed
in this subsection on a case-by-case basis. In the context of vectors, subtraction is to be
considered a case of addition, in which one is adding a vector whose direction is
somehow denoted as negative. The term for the sum of vector quantities is the
resultant.
In one dimension, the addition and subtraction of vector quantities is relatively
straightforward. If, for example, two vectors are parallel to each other, meaning they
have the same direction in one dimension, then the vectors are simply added together and
the resulting vector has the same overall direction.
Ex.) Consider the two vectors below
5u
3u
where u represents some unit (e.g. m/s). Since these vectors have the exact same
direction in space, they are simply added together, yielding a resultant with the same
direction as the original vectors:

104

8u
In a slightly more complicated case, the vectors are antiparallel, meaning they are
oriented in exactly opposite directions in one dimension. In this case, one must specify
which direction, left or right, is positive and which is negative. In accordance with the
conventional Cartesian coordinate plane, the right is usually the positive direction and the
left is usually the negative direction.
Ex.) Consider vectors of the same magnitude as before, only now oriented in opposite
directions:
5u
3u
If the left direction is considered the negative direction, then one is essentially adding
3 and -5. The resultant vector is:
2u
Notice how the vector of larger magnitude was more important in the sense that it
determined the direction of the resultant vector.
In more complicated cases, vectors lie in mores than one dimension. In the case of two
dimensions, one must apply laws of geometry and trigonometry. Furthermore, while
angles of either 0 (for parallel vectors) or 180 (for antiparallel vectors) were implied in
the examples concerning vectors in one dimension, angles must be explicitly indicated in
two dimensions. The following examples discuss common situations of vectors in two
dimensions.
Ex.) Suppose one were asked to determine the resultant vector of the following two
vectors:
4u
6u
How would one begin to solve this? First, one must exploit a fundamental property of
vectors that has already been discussed; vectors can be moved anywhere in space as long
as their magnitude and direction are preserved. Thus, to facilitate this problem, one
positions the vectors like so:

6u
4u

105

The resultant vector connects the tail of the smaller vector to the head of the larger
vector to form a right triangle:
6u

4u
The magnitude of the resultant can be found by using the Pythagorean theorem and the
angle at which this vector is displaced upwards from the horizontal can be found through
right-triangle trigonometry:
R 2 = (4u ) 2 + (6u ) 2 R = (4u ) 2 + (6u ) 2 = 2 13 u
6u
= tan 1 0.9827937232 56.30993247 o.
4u
What if the vectors cannot easily be situated 90 apart to form two sides of a right
triangle? One method often used exploits the key properties of a parallelogram; opposite
sides and opposite angles are equal. In this so-called parallelogram method, the two
vectors in question are situated next to each other in a tail-to-tail fashion such that they
form the two sides of a parallelogram. To find the resultant vector, which is the diagonal
of the parallelogram, one would use the law of sines and cosines.

Ex.) Consider two vectors, one that lies on a horizontal axis and one that is displaced
40 from the horizontal axis:
5u
40

3u
Since vectors can be moved anywhere in space as long as their magnitude and direction
are preserved, one can form a parallelogram from these two vectors by situating
themtail-to-tail:

40

Given that, in a parallelogram, opposite sides and opposite angles are equal, the complete
parallelogram has the properties depicted in the following diagram:

106

One can now use the law of cosines to determine the magnitude of the resultant vector
(i.e. the diagonal of the parallelogram):
20
3u
20
5u

140

A 2 = B 2 + C 2 2 BC cos a, letting A be the unknown resultant, B be 5u, C be 3u,

and a be 140:
R 2 = (5u ) 2 + (3u ) 2 2(5u )(3u ) cos(140 o ) R 7.548598101u.
Thus, the resultant vector has a magnitude of 7.548598101u and a direction of 20
from the positive horizontal axis.
In the last two problems, the resultant vector of two vectors in two-dimensional
space was found. How can one find the resultant vector of three or more vectors in twodimensional space? This task requires the process of vector resolution. In this process,
one breaks a vector down into its x- and y-components. Suppose the following vector lies
in the xy plane:

This vector is actually equivalent to the sum of two vectors, one that lies on the xaxis and one that lies on the y-axis. The magnitudes of these vectors are found by using
the following formulae:
Ax = A cos , Ay = A sin , where Ax and Ay represent the x- and ycomponents of vector A, repectively, A represents the magnitude of vector A (i.e., just
107

the number, which could also be represented by an unbolded A), and is the angle from
which vector A is displaced from the x-axis. Notice how the components of the vector
are not vectors themselves! This is because the formula from which they are derived
comes from right-triangle trigonometry, in which only magnitude is important. These
formulae should look familiar as they are essentially a conversion from polar coordinates
to Cartesian coordinates, as was discussed in the previous chapter. Vector resolution
represents both a mathematical and physical reality. Mathematically, resolving a vector
into its x- and y- components is basically an anti-sum, since one is finding the vectors
that could be added to yield the vector that one is resolving. This process also has
physical significance. Suppose someone is pulling a wagon with a force of 20 newtons
(N) at an angle of 60 from the horizontal. This person is actually pulling the wagon
forward and pulling the wagon up at the same time! Specifically, the wagon is pulled
forward with a force of (20 N) cos 60 o = 10 N and pulled upward with a force of
(20 N) sin 60 o 17.32050808. This person is actually pulling the wagon up to a greater
extent than he or she is moving it forward!
Vector resolution, as said before, is also useful in determining the resultant of three or
more vectors in two-dimensional space. In this process, one resolves each vector into its
x- and y-components (using the convention of measuring angles from the positive x-axis),
adds all of the x-components and all of the y-components, determines the magnitude of
the resultant vector ( R ) by converting from Cartesian coordinates to polar coordinates
2

with the formula R = R x + R y , and finally determines the angle at which this
resultant is displaced from the positive x-axis through the use of the formula
Ry
.
= tan 1
R
x
Ex.) Determine the resultant vector of the four vectors depicted below:

Start by resolving each vector into its x- and y-components, measuring all angles
from the positive x-axis:
A x = A cos = (10u ) cos 72 o = 3.090169944u
A y = A sin = (10u ) sin 72 o = 9.510565163u

B x = B cos = (12u ) cos 0 o = 12u

108

B y = B sin = (12u ) sin 0 o = 0u (Note that vector B has no y-component

because it lies completely on the x-axis)

C x = C cos = (15u ) cos 215 o = 12.28728066u (The negative sign
indicates that the x-component of vector C lies on the negative x-axis, but should not be
considered to indicate direction since components of a vector are not vectors themselves)
C y = C sin = (15u ) sin 215 o = 8.603646545u (Similarly, this ycomponent lies on the negative y-axis)
D x = D cos = (8u ) cos155 o = 7.250462296u
D y = D sin = (8u ) sin 155 o = 3.380946094u

One can now find the x- and y-components of the resultant vector, leading to the
magnitude of this vector:
R x = A x + B x + C x + D x = (3.09016994u ) + (12u ) + (12.28728066u ) + (7.250462296u )
= -4.447573016u
R y = A y + B y + C y + D y = (9.510565163u ) + (0u ) + (8.603646545u ) + (3.380946094u )
= 4.287864712u
R =

R x + R y = (4.447573016u ) 2 + (4.287864712u ) 2 = 6.177919514u (Notice how

the negative sign associated with the x-component R was not included because this
formula is derived from the Pythagorean theorem, and negative signs have no meanings
in classical geometry) The direction of R can be determined in the next step. Note that
the x-component of R will not be shown as negative here as well, since the following
formula is derived from right-triangle trigonometry, and triangles cannot have negative
dimensions:
Ry
4.287864712u
o
= tan 1
= tan 1
= 43.95258913 .
4.447573016u
Rx
Where exactly is this angle (i.e. in what quadrant)? To remove this ambiguity, one
examines the signs of Rx and Ry. Since the former was originally found to be negative
and the latter is positive, this angle lies in quadrant II. Thus, when measured from the
positive x-axis, the angle is 180 o 43.95258913o = 136.0474109 o. As a final answer, the
resultant vector is 6.177919514u, 136.0474109 from the positive x-axis.
This was quite a bit of information to deal with, so here is an overview of the theory
behind it:
Any vector can be resolved into its components, which are not truly vectors
themselves, due to their geometric and trigonometric derivation. Vector resolution is
particularly helpful in determining the vector sum (i.e. the resultant) of three or more
vectors in two-dimensional space. In this process, one first determines the x- and ycomponents of all the vectors in question. Note that these components do not
intrinsically have direction, though they may be positive or negative depending upon
their angle from the positive x-axis. One then sums all of the x-components and all of the
y-components, keeping the negative signs in the process. One then converts from

109

Cartesian coordinates to polar coordinates, using only positive values for Rx and Ry. Once
both the magnitude and associated angle of the resultant are found, one must determine in
which quadrant it lies. One does this by examining the original values of Rx and Ry.
There is a good intuitive way to informally check ones answer to a problem of this
type. Imagine that the vectors in question all represent forces tugging on something at
the origin. The resultant vector determines in which direction that thing will move.
Based upon the lengths and directions of the arrows representing the vectors, one should
be able to tell which magnitude and direction will win this strange game of tug-of-war.
In the examples already presented, the combined effects of vectors C and D are
qualitatively greater than those of A and B. Thus, the resultant vector will most likely be
pointed toward the negative x-axis, which was found to be true. Furthermore, the
combined effects of vectors A and D tugging upwards seems to overshadow the action
of vector C. One should, therefore, suspect that the resultant vector is pointed upwards,
which was also found to be true.
Unit Vector Notation and Generalization to Three Dimensions: There is a more concise
way to express the magnitude and direction of a vector than through the use of polar
coordinates. Indeed, an answer such as 6.177919514u, 136.0474109 from the positive
x-axis is rather burdensome when doing problems. Instead, one often wishes to express
a vector in terms of its components. However, as was discussed earlier, the components
of a vector are not vectors themselves, so simply displaying them does not suffice to
characterize a vector quantity. In order to use vector component notation and actually
describe a vector one must understand the concept of the unit vector. A unit vector is a
dimensionless (i.e. without units) vector that has a magnitude of unity (1 thus the name
unit vector) and that points along an axis. In two dimensions, a unit vector represented
by i points along the x-axis and a unit vector represented by j points along the y-axis:

How do unit vectors allow one to write a vector in terms of its components? By
multiplying these unit vectors to the components associated with the same axis, one can
easily express a vector by the sum of its components multiplied by these unit vectors. To
clarify, consider vector A below:

110

This vector has an x-component ax defined as a x = A cos and a y-component

defined as a y = A sin . If a x is multiplied by i and a y is multiplied by j, one can
represent vector A as follows:
A = Ax i + Ay j .
This is the representation of vector A in Cartesian coordinates. In polar coordinates,
Ay
2
2
this same vector is represented as A = Ax + Ay , tan 1 .
Ax
Note that when writing i and j by hand, the convention is usually to put a caret (^)
over the letter i or j: i and j . These are read as i hat and j hat, respectively.
Ex.) Express the vector below in both Cartesian and polar form.

Let the vector above be denoted as A.

Ax = A cos = (30u ) cos 50 o = 19.28362829u
Ay = A sin = (30u ) sin 50 o = 22.98133329u

Thus, in Cartesian form: A = (19.28362829u )i + (22.9813329u ) j.

In polar form, the vector is simply: A = 30u , 50 o from positive x-axis.
Note that on the AP Calculus BC exam, vector components might be expressed within
pointed brackets (< >). In this case, the values of the x- and y-components are put inside
the pointed brackets and separated with a comma. For instance, in the example above,
the AP exam might express it as A =< 19.28362829u, 22.9813329u > .
So far, vectors in only two dimensions have been discussed. While two-dimensional
systems often suffice to model real systems in nature, the universe is three-dimensional
(four if one includes time in the context of Einsteinian physics). Thus, it is often
necessary to work with vectors in three dimensions (often denoted as R 3 , meaning threespace). The formulae derived for vectors in two dimensions can be easily applied to
three dimensions by extending these formulae by one more element. Before discussing
the formulae in three dimensions, it is important to define an appropriate coordinate
system in which to work. Most textbooks use what is known as the right-handed
coordinate frame, thus named because when the fingers of the right hand are aligned
with the x-axis and curled toward the y-axis, the thumb represents the z-axis. The righthanded coordinate frame is shown below:

111

Notice that the extra dimension added is known as the z-axis. The unit vector
associated with this axis is denoted as k (or k if handwritten). Thus, one can now extend
the formulae for vectors in two dimensions to vectors in three dimensions:
A = Ax i + Ay j + Az k
A = A=

Ax + Ay + Az .

In the field of linear algebra, these formulae are extended to even more dimensions,
but that will not be discussed here.
Multiplication with Vectors: Before moving on to more difficult material in this
subsection, it is important to note some algebraic properties of vectors when
multiplication is concerned. If A and B are two vectors in three dimensions and m and n
are two scalars (i.e. with only magnitude, no direction), then the following algebraic
statements hold true:
mA = mAx i + mAy j + mAz k
(m + n) A = mA + nA = (mAx i + mAy j + mAz k ) + (nAx i + nAy j + nAz k )

n(mA) = m(nA) = mn Ax i + mn Ay j + mn Az k

m( A + B) = mA + mB = (mAx i + mAy j + mAz k ) + (mB x i + mB y j + mB z k )

The key point to take from these expressions is that when one multiplies a vector by
a scalar, nothing fancy is required; it suffices just to multiply every component of the
vector by the scalar. However, the same simplicity does not hold for multiplying a vector
by another vector. In general, there are two ways in which one can multiply two vectors.
One of these results in another scalar and one results in a vector.
The Scalar (Dot) Product: In general terms, the scalar product of two vectors in twodimensional space is described as the product of the magnitude of one vector and the xcomponent of the other vector. The scalar product is indicated by a dot () , hence the
alternative name dot product. The equation for the scalar product of two vectors A and B
separated by an angle is given as A B = A B cos = AB cos . The geometric
interpretation for this is shown below:

112

Thus, the scalar product is the product of the magnitude of a vector and the projection
of another vector onto that vector. While the name of this type of product certainly gives
it away, note that the result is a scalar quantity, not a vector quantity.
It is important to note special algebraic properties of the scalar product:
A B = Ax B x + Ay B y + Az B z This expression relates to vectors in R3. It basically
states that the scalar product is the sum of the ordinary products of the components in
each dimension.
A B = B A Scalar products are commutative.
A (B + C) = A B + A C Scalar products are distributive
m( A B) = (mA) B = (mB) A = m( A B) , where m is a scalar. This expression
extends the distributive law to scalars.
i i = j j = k k = 1, i j = i k = j k = 0. Recall that unit vectors have magnitudes of
unity. The first set of scalar products is equal to unity because they involve
multiplication of two vectors with the exact same direction (and magnitude). The second
set of scalar products is equal to zero because they involve multiplication by vectors that
are perpendicular. To see why this is the case, refer back to the diagram of the righthanded coordinate frame.
In physics, the most common application of the scalar product is the
determination of a quantity known as work (W), which is the scalar product of the force
(F) exerted on an object in newtons (N) and the displacement (s) (the change in position)
of that object in meters (m). Thus, W = F s = Fs cos . This equation only holds for a
constant force. Note that when the force and displacement vectors have the same
direction, the angle between them is 0. Thus, when the force is in the same exact
direction as the displacement, work is simply the product of F and s, since cos (0) = 1.
When the force vector is perpendicular to the displacement vector, no work is done on the
object in question because cos (90) = 0. Thus, while it may seem like a person would be
doing a great deal of work when lifting a cinder block while he or she is walking, this
person is doing absolutely no work in the physical sense!
Ex.) A wagon is pulled with a constant force of 25.0 N for a distance of 30 m as
shown by the diagram below:

113

Determine the work done on this object.

W = Fs cos = (25.0 N)(30.0 m) cos 45 o 539.3300859 539 Nm.
Note that the unit Nm (the newton-meter) is usually expressed by the equivalent unit of
the joule (J), which is also the unit for energy.
The Vector (Cross) Product: In contrast to the scalar product, the vector product yields a
result that is a vector quantity. Unfortunately, one cannot perform an operation that will
yield the magnitude and the direction of the result in one step. One first determines the
magnitude of the vector product and then applies a rule-of-thumb (literally!) to determine
its direction. The formula for the magnitude of the vector product between two vectors A
and B separated by an angle of is A B = A B = AB sin . Due to the notation in
this formula, the vector product is often referred to as the cross product. Before moving
on, it is important to note some important algebraic properties of the vector product.
A B B A, but A B = B A This means vector products are not commutative,
meaning the order of the vectors in the formula is very important. As will be discussed
shortly, this order determines the direction of the vector product.
A (B + C) = A B + A C Vector products are distributive.
m( A B) = (mA) B = (mB) A = m( A B) , where m is a scalar. This expression
extends the distributive law to scalars.
i i = j j = k k = 0, i j = k , j k = i, k i = j The first series of vector products
should make sense; since the angle between these unit vectors is 0, and since the sine of 0
is 0, the vector product is also 0. The second series of vector products should become
more clear in the discussion below.
One determines the direction of the vector product by using what is known as the
right hand rule. In this rule, one places the fingers of his or her right hand along the
direction of the first vector in the cross product and curls the fingers in the direction of
the second vector in the vector product. (Note: It must be the first vector because vector
products are not commutative, meaning, for instance, starting and B instead of A would
yield a completely different direction from what would result if one started at A) In
whichever direction the thumb points represents the same direction in which the vector
product points. Note that the vector product will always be perpendicular to the two
vectors in question. To understand this more clearly, observe the diagrams below:

114

This diagram represents the vector product A B . Notice that the fingers of the
right hand start at A and curl toward B. Since A came first in the equation, it cannot be
the other way around (i.e. one cannot start at B). The thumb points upward, so this is the
direction of the vector product A B .

This diagram represents the vector product B A . In this case, the fingers of the
right hand start at B and curl toward A. Since the thumb points downward, this is also
the direction of the vector product B A . Notice that this direction is the direct opposite
of A B . This supports the algebraic statement A B = B A . Note, however, that
both vector products are perpendicular to the plane in which the two vectors lie. In the
diagram, this was indicated by drawing in three dimensions. One could also use the
symbols and . The former indicates that the vector product is coming out of the
plane of the paper, while the latter indicates that the vector product is going into the
plane of the paper. Thus, the symbol would be associated with the first diagram and
the symbol would be associated with the second diagram. Note that there is a way to
determine the exact unit-vector notation of the vector product, but it involves a linearalgebraic tool known as the determinant, which will not be covered here.
Vector products are important in describing phenomena in many areas of physics.
Two applications will be discussed here. The first concerns rotational dynamics, or the
study of the forces that cause rotational motion. Everyone is probably familiar with the
fact that a door is more difficult to close when pushing closer to its hinges. However, one
can close the door almost effortlessly by exerting a force at the very end, away from the
hinges. Why is this the case? The force that one is exerting to cause rotational motion is
called the torque (represented by the Greek letter tau ). Torque is a vector quantity
that is directly proportional to the magnitude of the force applied and the distance from
the axis of rotation at which this force is applied. Thus, for a certain torque to be
achieved, if one starts at a small distance from the axis of rotation, a greater force would
need to be applied. Similarly, if one is farther away from the axis of rotation, one would
need to apply less force to achieve the required torque. Torque also depends upon the
angle at which the force is applied to the object in question. Applying the force

115

perpendicular to the object maximizes the torque while applying the force parallel to the
object causes no torque. Thus, a force that is applied perpendicular to a door will
maximize its rotational motion, while a force that is applied to the side of the door toward
the hinges will (obviously) cause no rotational motion at all. All of these elements can be
summarized by a vector product: = = F r = Fr sin , where F is the force applied,
r is the position relative to the axis of rotation at which the force is applied, and is the
angle between the force and position vectors.
Ex.) The following diagram represents a 3.00-m long wooden rod attached to a hinge
that allows the rod to rotate freely on its axis with negligible friction.

If a force of 50.0 N is applied to the wooden rod at an angle of 30 upward from

the left, at what distance from the axis of rotation must this force be applied to generate a
torque of 6.75 Nm? (Note that the units of torque are newton-meters, but unlike work,
are not considered to have units of joules).
In this problem, it would be beneficial to draw out the force and position vectors:

Unfortunately, this is not in a form conducive to the vector product. In the vector
product, the vectors are connected head-to-head like so:

A little geometric analysis would show that, luckily, the angle between these two
vectors in this situation is also 30. With this information, one can find the magnitude of
the position vector:

116

(6.75 N m)
= 0.27 m .
F sin (50.0 N) sin 30 o
Note that while it was not explicitly a part of this problem, the direction of the torque
would be into the page ( ) because the fingers of the right hand are starting at the
force vector (because it appears first in the equation) and curling toward the position
vector.
Another interesting application of the vector product is in the field of
electromagnetism. Electricity and magnetism, due to the work of the 19th-century
Scottish physicist James Clerk Maxwell, are united based upon the fact that electric
charges in motion generate magnetic fields. This is fundamentally true, even in, say, a
refrigerator magnet! It is also the case that a preexisting magnetic field will exert a force
upon a moving point of charge. This force is directly proportional to the magnitude of
the charge, the velocity of the charge, and the strength of the magnetic field. The force
exerted on the point of charge by the external magnetic field also depends on the angle at
which the charged particle moves relative to the magnetic field. The magnetic force is
maximized when the charge moves perpendicular to the magnetic field and is zero when
it is parallel to the magnetic field. All of these observations should signal that the force
exerted by a magnetic field on a charged particle can be represented by a vector product:
F = F = qv B = qvB sin , where F is the force of the magnetic field, q is the

= Fr sin r =

magnitude of the charge (Note that charge is always a scalar quantity), v is the velocity of
the charged particle, B is the strength of the magnetic field, and is the angle between
the velocity and magnetic fields vectors. Charge magnitude is measured in coulombs (C)
and magnetic field strength is measured in newtons per coulomb-meters per second
(N/Cm/s), which is succinctly referred to as the tesla (T).
Ex.) Consider a charged particle that is moving out of the page in a magnetic field that
is directed to the left:

This particle has a charge of 9.1314 10 19 C and is moving out of the page with a
velocity of 3.3017 10 2 m/s . If the strength of the magnetic field is 2.7877 T and the
charged particle experiences a force of 5.1860 10 16 N, what is the angle between the
velocity vector and the magnetic field vector?

F = qvB sin sin =

F
(6.1860 10 16 N)
F

= sin 1
= sin 1
19
2
qvB
qvB
(9.1314 10 C)(3.3017 10 m/s)(2.7877 T)

= 47.393o.

117

Vector-Valued Functions
Note: This section is on the AP Calculus BC exam!
Thus far, the theory and application of vectors have been extensively discussed. This
section introduces perhaps an even more highly powerful use of vectors in mathematical
analyses. Scientists and mathematicians often express functions in terms of vectors to
analyze curves in space, and this will be the theme for the remainder of the chapter. A
vector-valued function, or simply vector function, is, not surprisingly, a function defined
by vectors. More formally, if real scalar values denoted by u can be located with a vector
A, then A is function of the scalar u: A(u). Suppose that the scalar values have x-, y -,
and z-coordinates. If this is the case, then A(u) can be expressed in terms of the
component vectors of A: A(u ) = Ax (u )i + Ay (u ) j + Az (u )k. This is the point where the
relationship between vector functions and parametric equations comes into play; each
component vector of the vector function is, itself, a parametric equation with u as the
parameter! To see why this is the case, consider a more familiar example from physics
the position function. This is actually a vector function denoted as r. A particles
position in three-dimensional space (i.e. x-, y -, and z-coordinates) with respect to the
time parameter t can be found through the vector function r (t ) = x(t )i + y (t ) j + z (t )k .
This can also be understood pictorially:

This diagram depicts the position vector r locating a particle in R 3 at a certain point
(x, y, z). It is important to note that this particles path is actually a parametrized curve.
The parametric equations are x(t ), y (t ), and z (t ). Thus, a vector function is essentially a
sum of parametric equations with the addition of unit vector notation. As stated in the
introduction to this chapter, this is basically the extent to which the AP exam tests the
theory of vectors.
The Calculus of Vector-Valued Functions
Note: This section is also on the AP Calculus BC exam! Note also: Now that calculus is
being applied again, be sure the graphing calculator is in radian mode!
The theorems of calculus as applied to vector functions are exactly analogous to those
of functions not defined by vectors. The only fundamental difference is that the calculus
theorems of vector functions include vector notation. For example, the limit definition of

118

the derivative is not necessarily unique for vector functions. For a vector function A(u ),
the limit definition of the derivative is as follows:
dA
A(u + u ) A(u )
= lim
. This is exactly analogous to the definition
du u 0
u

introduced at the beginning of AP Calculus AB. Furthermore, the rules for

differentiating vector functions are also analogous to those already learned. For two
vector functions A(u ) and B(u ),
d
dA dB
( A + B) =
+
du
du du
d
dA
dm
(Product rule for a vector and a scalar m)
(mA) = m
+A
du
du
du
d
dB
dA
( A B) = A
+B
(Product rule for a scalar product)
du
du
du
d
dB
dA
( A B) = A
+ B
(Product rule for a vector product).
du
du
du
In most calculus courses, including AP Calculus BC, one applies vector differentiation
most often to the equations of motion position, velocity, and acceleration. Little has
changed, however, in regards to the definitions of these quantities; velocity is still the
time-derivative of position and acceleration is still the time-derivative of velocity. It is
just that, now, one must recognize these equations in vector notation. Beginning with the
position function r (t ) = x(t )i + y (t ) j + z (t )k , one can derive the expressions for velocity
and acceleration:
dr (t ) d
d
d
= x(t )i + y (t ) j + z (t )k = v x (t )i + v y (t ) j + v z (t )k
dt
dt
dt
dt
dv (t ) d
d
d
a(t ) =
= v x (t )i + v y (t ) j + v z k = a x (t )i + a y (t ) j + a z (t )k
dt
dt
dt
dt
v (t ) =

It is also important to know the following formulae in relation to particle motion in vector
space:
2

speed = v (t ) = v x + v y + v z

a x (t ) =

Speed, a scalar quantity, is the magnitude of velocity

d v(t )

The horizontal component of the acceleration vector yields the change

dt
in speed of the particle. This type of acceleration is also referred to as tangential
acceleration.
a y (t ) describes the change in direction of the particle. This type of acceleration is
also referred to as radial acceleration. Thus, each component of acceleration describes a
different aspect of the velocity vector; the x-component of acceleration describes the
change in speed (i.e. the magnitude) and the y-component describes the change in
direction.

119

Integration of vector functions is also quite straightforward, for the rules of integration
are the same. Integrating a vector function means integrating all of the components of the
vector function. For a vector function A(u ),

A(u)du = A (u)i du + A (u) j du + A (u)k du .

x

Essentially, integrating a vector function merely involves summing the integrals of

parametric equations! Again, the AP Calculus BC exam tests extensively on this point.
Thus, to do problems of this sort, it may very well suffice to understand parametric
equations and vector notation without understanding vectors at all! I do not recommend
this, however; otherwise, I would not have introduced their theory for the first 14 pages!
Similar to differentiation of vector functions, integration of vector functions is usually
applied to the equation of motion. Again, nothing fundamental has changed; velocity is
still the integral of acceleration and position is still the integral of velocity. However, one
must now understand these equations in vector notation:

a(t )dt = v(t ) + C, where the constant of integration C is also a vector:

C = Cxi + C y j + Czk .

v(t )dt = r(t ) + C

Given the appropriate initial conditions, one can determine vector C. It is also possible
to perform a definite integral to find, for example, the change in position of an object
over a certain time period:

t2

v (t )dt =

t1

t2

v x (t )dt +

t1

t2
t1

v y (t )dt +

t2
t1

v y (t )dt = r

Is there significance for the definite integral of the speed of a particle from one time to
another? To see if there is, one can examine the integral and see if it looks like anything
familiar. Suppose one is working in two dimensions rather than three. The definite
integral of speed would have the following appearance:

t2
t1

v (t ) dt =

t2
t1

v x + v y dt =

t2
t1

dy
dx
+ dt.
dt
dt

This last integral should look familiar; it is the formula for the arc length of a
parametric curve! Since this is the length of the particles path, the definite integral of
the speed of a particle from one time to another equals the total distance distance traveled
by the particle during that period of time. What is the difference between change in
position (displacement) and total distance traveled? Consider a particle moving along the
following path in two dimensions:

120

The change in position (displacement) of the particle is r = ( x 2 i + y 2 j) ( x1i + y1 j) .

This is merely the difference between the two indicated points. Notice that change in
position is a vector quantity. The total distance traveled, on the other hand, is the
complete length of the curve from the initial point to the end point. The total distance
traveled is a scalar quantity.
Problems concerning the different aspects of the motion of a particle with vector
notation included appear very frequently on the free response portion on the AP Calculus
BC exam. I stress again that these problems are really just a test of knowledge of
parametric equations in relation to particle motion.
Ex.) A particle is moving in the xy plane with a velocity given by the vector-valued
function v(t ) = (5t e t )i + (3t 3 )j. Initially, the particle is located at the origin.
a.) Derive the position function for this particle.
b.) Determine the speed of the particle at t = 3 to three decimal places.
c.) Find the total distance traveled by the particle from t = 0 to t = 4 to three
decimal places.
d.) Find the equation of the line tangent to the particles path at t = 3.
a.) This part requires integration of the velocity function and knowledge of the initial
conditions, which are given:
r (t ) =

v (t )dt =

(5t e )idt + (3t )jdt

t

The integral of the x-component of velocity requires integration by parts, but since
one term is a polynomial, one can use the tabular method (see chapter 1):
u
dv
1
t
5t
+1
e
-1
5
e t

t
+1
0
e
-1
Thus,

(5t e

)idt =( 5t e

5e t )i + C x = x(t )i + C x . Since the x-coordinate of the

particle at t = 0 is 0, Cx is also 0, so x(t )i = (5t e t 5e t )i .

121

To find the y-component of the position, one must integrate the y-component of the
velocity:
3
3t 3 j dt = t 4 j + C y = y (t ) j + C y . Since the y-ccordinate of the particle at
4
3
t = 0 is 0, Cy is also 0, so y (t ) j = t 4 j . Thus, the position equation is:
4
3
r (t ) = 5t e t 5e t i + t 4 j.
4
b.) Recall that speed is the magnitude of the velocity vector. So, at t = 3:

speed t =3 = v (3) =

[5(3)e ] + [3(3) ]
( 3) 2

3 2

= 6561.558.

c.) The total distance traveled is the definite integral of speed, or the arc length of the
parametric curve, so:
total distance t =0,t =4 =

v x 2 + v y 2 dt =

(5t e ) + (3t )dt.

4

t 2

Using the graphing

calculator, total distance = 7.770.

d.) With the information given, one can use the point slope formula
( y y1 = m( x x1 )) to supply the equation for the tangent line at t = 3. The slope if the
dy
at t = 3 and the x- and y-coordinates can be found from the position
value of
dx
equation:
dy

dy
3(3) 3
= dt =
= 108.4618994.
dx t =3 dx
5(3)e ( 3)

dt t =3
x(3) = 0.9957414
y (3) = 60.75
y (60.75) = (108.4618994)( x + 0.9957414) .
The Differential Geometry of Curves in Space

Differential geometry is a very rich field of higher mathematics. Generally

speaking, the study of differential geometry applies the methods of calculus to geometric
entities like curves and surfaces. This field is of particular importance not only in pure
and applied mathematics, but in physics (especially general relativity) and theoretical
chemistry (especially molecular modeling). Of course, this section cannot possibly do
justice to the entire field of differential geometry. Rather, here will be discussed the
analyses of various aspects of curves in space. For these analyses, it is necessary to work
with vectors.

122

The TNB Frame: In the study of curves in space, a generic vector function is usually
represented as the position function: r (t ) = x(t )i + y (t ) j + z (t )k. This function is
assumed to be differentiable. The material introduced in this subsection is most
popularly applied to motion of particles in space, especially as applied to flight.
Nevertheless, the position function is still useful in analyzing curves that do not model
motion. For instance, the double-helix model of a DNA molecule can be analyzed from
the standpoint of the position function as well.
In R 3 , curves are often analyzed through the use of three types of vectors. The
first is the unit tangent vector (T). This vector is useful because it quantifies the
direction in which a particle is moving at a certain point on a curve. A special property
of the unit tangent vector is that it is magnitude-blind; it provides a measurement solely
v
for direction. This is evident in the formula for the unit tangent vector: T = . Notice
v
how division by the magnitude of the velocity vector (speed) allows the unit tangent
vector to solely quantify direction. Note also that this division results in a magnitude of
unity for this vector. The unit tangent vector is also sometimes conveyed as the
derivative of the position function with respect to arc length (which, in this case, is
dr
represented by s). Thus, it is also the case that T = . This formula is very impractical
ds
because it is quite difficult to find position in terms of arc length. The second vector used
to analyze curves in space is called the principal unit normal vector (N). This vector is
orthogonal (i.e. perpendicular) to the unit tangent vector. The principle unit normal
vector primarily conveys the extent to which a particles path deviates from a straight
line. Its formula is derived from a scalar quantity known as curvature ( the Greek
letter kappa). Curvature (besides the obvious definition of extent to which a curve
curves!) is related to the rate at which the curve changes direction at a point. Thus, it is
defined mathematically as the magnitude of the derivative of the unit tangent vector with
dT
respect to arc length: =
. Again, it is impractical to manipulate a function to be
ds
defined by arc length. Thus, one can use the chain rule in the following way:
dT ds
=
, where the latter term is simply the reciprocal of speed. It is also possible to
dt dt
simplify the expression of even more by expressing it in terms of familiar aspects of
motion. While the derivation will not be shown here, curvature can also be found from
va
the following formula: =
(Note that the numerator is a vector product). How does
3
v
curvature relate to the principle unit normal vector? By multiplying the derivative of the
unit tangent vector with respect to arc length by the reciprocal of curvature, one yields the
dT 1
principle unit normal vector: N =
. Unfortunately, this equation is in terms of arc
ds
length, so the chain rule is used to yield a formula in terms of t:

123

dT dt
dT
1
d T 1 dT
=
= dt ds = dt . Due to this division, the principal unit
N=
dT dt
dT
ds ds dT / ds
dt ds
dt
normal vector also has a magnitude of unity. The third vector used in curve analysis is
the unit binormal vector (B). This vector is orthogonal to both the unit tangent vector
and the principal unit normal vector. Loosely speaking, the unit binormal vector
measures the extent to which a moving particle is dissociating itself from its path. The
unit binormal vector is simply the vector product of the previous two vectors discussed:
B = TN .
It is now time to put all of this information together. The unit tangent vector, principle
unit normal vector, and unit binormal vector all have unit length and lie orthogonal to
each other. Thus, they form a three-dimensional coordinate frame, known as the TNB
frame, or Frenet frame, in honor of the 19th-century French mathematician Jean Frdric
Frenet (pronounced fruh-nay):

This frame is always associated with the particle in question. That is, as the particle
moves along a curve in space, the TNB frame always moves along with it. Throughout
the particles journey, the TNB frame describes different aspects of that particles
motion; T describes the direction of the particle at a particular point, N describes the
extent to which that particles path deviates from a straight line, and B describes the
extent to which a particle could potentially take another path besides the one it is
actually taking. It may be more helpful to understand these aspects of particle motion
from a common experience (for some) riding on a roller coaster. In this case, T would
represent the direction tangent to the coaster one feels that he or she is moving, N would
represent the extent to which one feels like he or she is curving, and B would represent
the extent to which one feels as if he or she is being pushed up. The following diagram
represents a particle traveling along a curve in space and the TNB frame associated with
it at specific points on the curve:

124

There is one more useful aspect of particle motion on curves in space to be discussed
torsion ( ). Torsion is a scalar quantity that describes the extent to which a curve
twists in space. Torsion is equal to the magnitude of the rate at which the direction of
dB
the unit binormal vector changes with respect to arc length: =
. There is a way to
ds
simplify this formula, but it involves the determinant, which is beyond the scope of this
book..
As one can glean from the formulae in this subsection, problems involving the TNB
frame are lethally tedious, and it is almost always the case that even the most
computationally careful of humans will make some sort of error. Thus, in the upcoming
example, the calculations will be done on Mathematica.
Ex.) A malfunctioning model airplane moves through space as described by the position
equation: r (t ) = sin(3.33t )i + cos(2.22t ) j + 1.11(t )k.
a.) Determine the curvature of the airplanes path at three times:
t = 0, t = , and t = 2 .
b.) Compare the equation of the derivative of the unit tangent vector with respect to
time to the equation for curvature. What is being compared?
It would be helpful to visualize this curves path using Mathematica :
ParametricPlot3D[{Sin[3.33*t],Cos[2.22t],1.11*t},{t,0,2*Pi}]

-0.5
5
-1

0.5
0

-1
1 -0.5 0 0.5

Graphics3D

Indeed, this plane is experiencing some problems. a.) To compute the curvature at the
three desired times, one uses the [r_][t_] function in Mathematica :
[gamma][0] = 0.4,[gamma][Pi]= 1.52186,[gamma][2*Pi]=0.897866

125

Note that [gamma]merely stands for the position function, for the purpose of
saving space. This plane is not even screwing up consistently! Its curvature increases
and then decreases on the time interval.
b.) For this part of the problem, one can use the same Mathematica operator,
but leave the result in terms of t. For instance, the curvature function can be found in the
following way:
,
2
2
I H29.9266 Cos@2.22 tD + 151.504 Sin@3.33 tD +

[gamma][t]=

1. H16.4116 Cos@2.22 tD Cos@3.33 tD +

24.6174 Sin@2.22 tD Sin@3.33 tDL2LM

H1.2321 + 11.0889 Cos@3.33 tD + 4.9284 Sin@2.22 tD L

2 32

In conventional mathematical terms, this expression would look like

29.9266 cos(2.22t ) 2 + 151.504 sin(3.33t ) 2 + 1.[(16.4116 cos(2.22t ) cos(3.33t ) + 24.6174 sin(2.22t ) sin(3.33t )]

[r (t )] =

[1.2321 + 11.0889 cos(3.33t ) 2 + 4.9284 sin( 2.22t ) 2 ] 3

Based upon the thoroughly confusing algebraic structure of this expression, it is

clear that it would not be practical to do such a problem by hand!
A similar computation can be done on Mathematica to determine the equation
for the unit tangent vector:
::

D[{UnitTangent[gamma][t]},t]
11.0889 Sin@3.33 tD

!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
1.2321 + 11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2

H1.665 Cos@3.33 tD H21.8821 Cos@2.22 tD Sin@2.22 tD

73.8521 Cos@3.33 tD Sin@3.33 tDLL H1.2321 +

11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2L

4.9284 Cos@2.22 tD

32

!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!!
1.2321 + 11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2

H1.11 Sin@2.22 tD H21.8821 Cos@2.22 tD Sin@2.22 tD

73.8521 Cos@3.33 tD Sin@3.33 tDLL H1.2321 +

11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2L ,

H0.555 H21.8821 Cos@2.22 tD Sin@2.22 tD
73.8521 Cos@3.33 tD Sin@3.33 tDLL H1.2321 +
32

11.0889 Cos@3.33 tD2 + 4.9284 Sin@2.22 tD2L

>>

32

Notice how this expression, which is even more confusing, gives three
coordinates instead of a single equation. This is the case because this is the derivative of
the unit tangent vector with respect to time, while the curvature was the derivative of the
unit tangent vector with respect to arc length.
126

The Calculus of Scalar and Vector Fields

This final section of the chapter very strongly converges on material characteristic of
Calculus III. I believe that this section provides an appropriate bridge into the study of
multivariable calculus.
Thus far, this chapter has largely been concerned with vector functions with respect to
time. However, it is often necessary to express quantities in terms of changes in space.
For instance, a meteorologist might express the variation is pressure or temperature over
a geographic region without any direct regard for time. A region of space with some
physical quantity assigned at every point is known as a field. These physical quantities
may either be scalars or vectors. Thus, one may describe either a scalar field, a field that
assigns scalar quantities at every point in space in a region, or a vector field, a field that
assigns vector quantities at every point in space in a region. Examples of scalar fields
include the electrical potential around a charged particle and the variation in certain
meteorological quantities in a geographic region. Examples of vector fields include, the
gravitational or electric force as it varies in space and flux, which describes the rate of
flow of some quantity at various regions in space. While fields can be represented
pictorially, such as the isotherms that connect regions of space with the same temperature
as shown on a weather map, recall that a field has a value assigned to every point in
space. To determine these values, one requires the continuity of a function. These
functions, known as field functions, mathematically express how a certain quantity
varies in space. If this quantity is a scalar, the corresponding function is called a scalar
field function. If it is a vector, the corresponding function is called a vector field
function. For example, the function F( x, y, z ) expresses force as a function of threedimensional space. It is important to notice that while previous sections concerned vector
functions as a function of time, a scalar quantity, this section concerns vector functions as
functions of position in space, which is a vector quantity. Thus, the calculus associated
with field functions is considerably different from that associated with time-dependent
vector functions.
Differentiation in Scalar Fields: One of the major differences between differentiation of
time-dependent and time-independent vector functions is that the latter requires
multivariable calculus since the function is of three variables of space x, y, and z. First
introduced in chapter 5 of this book, one differentiates a function of multiple variables by
using a partial derivative, which is basically a derivative with respect to one element of
the function that is allowed to vary while keeping the other elements constant. For
instance, the partial derivative of T ( x, y, z ) with respect to x while keeping y and z
T
. When differentiating field functions, one is
constant would be written as
x y , z
concerned with spatial variation of some quantity. In the context of scalar quantities,
there are two primary measurements used to examine the spatial variation in those
quantities. The first is the directional derivative. This derivative expresses the change
in some scalar quantity from one point in space to another. Manifested in its name is its
ability to convey a change in a specific direction. Thus, the directional derivative is a
vector quantity. As an example, consider a region of space in which the scalar quantity
127

assigned to each spatial point is temperature in degrees Celsius. Assume that in the
region being studied, temperature only varies in the x-direction:

If T is a function of x, y, and z and temperature only varies in the x-direction from T1 and
the vector u indicates the direction from T1 to T2 then the directional derivative
suggested by the diagram is:
T T
T
u = 2 1 u.
x y , z
x 2 x1
While the directional derivative in this case was chosen to be from T1 to T2, it could have
been from T1 to any other temperature. Thus, given a point in space that represents a
scalar quantity, there are an infinite number of directional derivatives from that point.
One directional derivative is especially significant. It is called the gradient. The gradient
is a vector that points in the direction of greatest increase in the scalar quantity in
question. The gradient is often denoted by a symbol, which is referred to as the del
operator. The gradient of a scalar field function f is defined mathematically as:
f
f
f
f =
i+
j + k.
x
z
dy
Gradients appear extensively in the physical sciences. For instance, the one-dimensional
dU
equation that relates force ( Fx ) to potential energy (U) is Fx =
. This can be
dx
U U U
+
+
.
extended to three dimensions through the use of the gradient: F = U =
x
z
y
In electrostatics, the gradient relates the electric field (E), a vector, to electrical potential,
V V V
or voltage (V), a scalar: E = V =
+
+
.
x y z
Differentiation in Vector Fields: Differentiation in scalar fields requires at most three
partial derivatives, one for each space component. Differentiation in vector fields,
however, requires more partial derivatives because both elements of the derivative are
vectors, each with different components. Consider differentiation in a force field
F( x, y, z ) . Unlike differentiation in a scalar field of temperature, the force vector itself
has three components; Fx , Fy , and Fz ; each of which are a function of x, y, and z. Thus,
for differentiation in a three-dimensional vector field, there are a total of nine different
partial derivatives! It is quite impractical to describe a physical system in this manner.
Thus, two new mathematical expressions are introduced that more concisely describe
128

spatial variation in a vector field. The first is divergence, which, for a vector field
function F, is often denoted as divF . Physically speaking, divergence measures the
extent to which a point in space is a source or sink for something, making it an ideal
quantity in physical systems characterized by the flow of matter and energy.
Mathematically, divergence is the scalar product between the del operator and the vector
field in question. For a vector field F:
Fy Fz
F

divF = F = i +
j + j ( Fx i + Fy j + Fz k ) = x +
. Note that
+
y
z
x
y
z
x
divergence is a scalar quantity. The second mathematical expression often used to
describe spatial rates of change in vector fields is called curl, which, for a vector field
function F, is often denoted as curl F. Note that curl is bolded because it, unlike
divergence, is a vector quantity. In physical terms, curl can be described as the extent of
rotation of a vector field, making it an ideal mathematical expression to use in fluid
dynamics, such as studying vortices. While divergence is the scalar product between the
del operator and the vector field in question, curl is the vector product between these two
elements (The last part of this expression was arrived at through the determinant):
F Fy Fx Fz Fy Fx

i +
k.
curl F = F = i +
j + j ( Fx i + Fy j + Fz k ) = z

j +
y
z

y
z
z
x
x
y

The Line Integral: This last subsection introduces an important mathematical expression
in vector calculus the line integral. This kind of integral, also sometimes referred to as
a path integral, measures scalar or vector quantities along a curve in space. One of the
most popular examples of a line integral concerns work, a physical quantity that was
introduced in the first discussion of the scalar product in this chapter. This section,
however, will consider the work done by a non-constant force. Consider a particle
moving along a path C described by the position function r. This particle is acted upon
by a non-constant force F(r) along this path. Recall that the work done on an object is
the scalar product of force and displacement. This formula still holds here, but since a
non-constant force is doing the work, one must move to the differential level and divide
the particles path into an infinitely large number of infinitesimal pieces in which the
force is constant. Taking the sum of all of these instances of work yields the actual
work done by the non-constant force. This, of course, involves an integral, specifically a
line integral: W =

F(r) dr. Note the new notation of this integral. It is not a definite
C

integral, but the subscript C signals that this integral is being evaluated over a curve. In
order to evaluate such an integral, it is necessary to manipulate it to yield an equivalent
definite integral. In realizing that the vector function r is actually a parametric equation
with the time parameter t, one can also express the force equation in terms of t:
F(r ) = Fx (t )i + Fy (t ) j + Fz (t )k. Taking the scalar product of the parametrized force and
differential position functions yields the following:
dy
dx
dz

+ Fz (t ) dt. Integrating this expression from

F(r ) dr = Fx (t ) + Fy (t )
dt
dt
dt

t1 to t 2 is equivalent to the line integral already introduced:

129

W =

F(r ) dr =

t2

t1

dy
dz
dx

Fx (t ) + F y (t ) + Fz (t ) dt .
dt
dt
dt

While this line integral involved motion over a so-called open path, line integrals can
also be evaluated over closed paths, in which the initial point coincides with the endpoint.
For instance, consider the closed path below:

If a non-constant force is acting upon a particle traveling along this curve, the notation
of the line integral must indicate that the path is closed. This is done by the inclusion of a
small circle around the integral symbol:

. The line integral describing the work done

on the particle can be expressed as the sum of the line integrals from point a to b to c and
point c to d to a:
W =

F(r ) dr =

abc

F(r ) dr + F(r ) dr.

cda

Oftentimes, a line integral only depends upon the points along the curve being
considered, not on the actual path taken. For instance, consider three possible paths from
point A to point B below:

While each path is rather distinct from the others, it may very well be the case that the
line integral is the same for each path. In this case of path independence, only the initial
and final points have any importance in the value of the line integral, not the actual path
taken. If this is the case, the field in question is called a conservative field. For instance,
the work done by gravity is irrespective of the path of an object,t while the work done by
friction does depend upon the path. Thus, gravity is an example of a conservative force
and friction is an example of a non-conservative force.

130

Concluding Remarks
This final chapter introduced the theory and calculus of vectors, a very important
concept in mathematics and science. The nature of vectors and the algebra that reflects
this nature were discussed. After vectors were considered from this relatively unique
standpoint, they were then discussed in the context of functions, which conveyed their
close relationship with the parametric coordinate system, especially in terms of position
velocity and acceleration. It was discussed that the theorems of vector calculus are
essentially the same as those for scalar functions, and this calculus was applied to
problems in motion. The final two sections introduced concepts that are generally
introduced in Calculus III. The penultimate section introduced the theory of the
differential geometry of curves and the mathematical expressions used to analyze curves
in space. The final section heavily converged upon multivariable calculus in discussing
the calculus of scalar and vector fields.
Key Terms:
vector quantities
scalar quantities
vector
resultant
vector resolution
unit vector
right-handed coordinate frame
scalar product
work
vector product
right hand rule
torque
vector-valued function
differential geometry
unit tangent vector
principal unit normal vector
curvature
unit binormal vector

TNB frame
torsion
field
scalar field
vector field
field functions
scalar field function
vector field function
directional derivative
gradient
del operator
divergence
curl
line integral
conservative field

131

Practice AP Calculus BC exam

I have created a practice exam based upon research of past multiple choice and free
response questions. I have tried to make this practice exam similar to what one may
actually encounter on the actual one. Most of these questions are in the style of actual
exam questions that I have kept from my AP Calculus BC class in high school. Most of
the questions on this practice exam are slightly more difficult than those that one would
encounter on the actual exam. I feel that this scenario is always better so that one does
not become too comfortable before the exam, only to become unpleasantly surprised on
exam day.
As stated in the introduction to this book, a good number of questions on the AP
Calculus BC exam test on material from AP Calculus AB. While I have provided
explained answers to the questions in the following practice exam, more extensive
discussions of the AB material do not appear in this book, which is specifically
concentrated on BC material. For such discussions, one should consult a review book
geared toward AB material.
Good luck!

132

Practice Exam
Section 1, Part A
55 minutes
28 questions
Directions: Solve each of the problems in this section, using the available space for
scratchwork. No credit will be awarded for anything written in the test booklet. Use
your time wisely. No calculator of any kind may be used in this section.
Note: Unless otherwise indicated, the domain of a function f is the set of all real
numbers x for which f (x) is a real number.
1.) If 3xy = x 4 y 3 , then

a.)

dy
=
dx

4x 3 y 3 3 y
3y 4 y3 x3
3x 4 y 2 + 4 x 3 y 3
, c.)
, b.) 4 2
, d.)
3
3x 3x 4 y 2
x 3y 3

9x 4
2x8

3x
x4

, e.)

9x 4
2 yx 8

________________________________________________________________________
2.) The growth of a population is modeled by the differential equation
dN N
N
= 1 , where N is the number of individuals in the population. How many
dt 12 72
individuals will there be when the carrying capacity is reached?
72

a.) 12, b.) 72, c.) e 12 , d.) ln (72), e.) ln (12)

________________________________________________________________________

133

3.)

3 sin xdx =

1 1
1
cos t , d.) 3 cos t , e.) cos t
3 3
3

a.) 3 3 cos t , b.) 3 cos t 3 , c.)

________________________________________________________________________
4.) Consider the area enclosed by the two curves shown below:

The formula for the volume of the solid generated by revolving the shaded area about
the y-axis is
5. 5

b.) 2 ( g ( x) f ( x) )xdx

c.)
[1 (g ( x) f ( x))] xdx
d.)
[1 + (g ( x) f ( x))] ydy
e.) 2
(g ( x) f ( x))ydy

a.)

(g ( x) f ( x) )2 dx

5. 5

5. 5

f (5.5)

f (0)

f ( 5. 5 )

f (0)

134

________________________________________________________________________
5.) A particle moves along a path described by the parametric equations
x(t ) = 5 sin t and y (t ) = t 3 . At t = 4, the particle leaves the path and begins to travel along
the tangent line to the path at that point. What is the slope of this line?
a.)

3 2
5
5 cos t
5
3
, e.) t 2 sec t
t sec t , b.) , c.) , d.)
2
5
3
3
5
3t

________________________________________________________________________
6.)

a.)
b.)
c.)
d.)
e.)

2x3 + 2x 2 + 1
dx =
x 4 + 2x3 + x

1
ln x 4 + 2 x 3 + x + C
6
1
ln x 4 + 2 x 3 + x + C
2
3
2 ln x + ln x 3 + 2 x 2 + 1 + C
4
1
ln x + 3x 3 + 2 x + C
2
1
1
2 ln x + x 3 + x + C
3
2

________________________________________________________________________
7.) Which of the following series converge?

I.)

n =1

n2 + 3
, II.)
n2

n =1

2n!
, III.)
10 n

n =1

3n
4n + 1

a.) I only, b.) II only, c.) III only, d.) Both I and III, e.) Both II and III

135

________________________________________________________________________
8.) If y = 3 ln(3 x) + 5 , then y =
a.)
b.)
c.)

1
3 3 [ln(3x) + 5] 2
1
3( x + 5) 3 [ln(3x) + 5] 2
1

3x 3 [ln(3 x) + 5] 2

d.) 3x 3 [ln(3x) + 5] 2
e.) 3 3 [ln(3x) + 5] 2

________________________________________________________________________
9.)
a.)
b.)
c.)
d.)
e.)

3 x 3 cos x dx =

3 4
x sin x + C
4
3x 3 sin x + 9 x 2 cos x + C
3x 3 sin x + 9 x 2 cos x 18 x sin x 18 cos x + C
3x 3 sin x 9 x 2 18 x sin c 18 cos x + C
3 4
3 5
3 6
3 7
x cos x +
x sin x
x cos x
x sin x + C
4
20
120
840

136

________________________________________________________________________
10.) The function f (x) is defined piecewise as:
2
f ( x) = x 1, x < 2
3a 2, x 2

For what value of a will the function f (x) be continuous?

a.)

1
3
1
5
3
, b.) , c.) , d.) , e.)
2
4
3
3
5

________________________________________________________________________
11.) The graph below represents f (x) , a functions derivative.

On the x-interval shown, what is the total number of minima and maxima on the graph
of the function f (x) ?
a.) 5, b.) 6, c.) 7, d.) 8, e.) Not enough information provided.

________________________________________________________________________

137

1
12.) lim 2 =

a.) 1, b.) 0, c.) , d.) e, e.) e2

________________________________________________________________________
13.) Find the value of x that satisfies the Mean Value Theorem for f ( x) = x 3 + 2 x + 1 on
the closed interval [0, 2].
a.) 1, b.) 6, c.)

2
3

, d.)

3
, e.)
2

________________________________________________________________________

138

14.)

a.)

dx
4 + 9x 2

1
3
3
1

tan 1 , b.) tan 1 , c.) , d.) , e.) divergent

12 6
8 2
2
2
6
8

________________________________________________________________________

15.) Given that the power series for e is

n =0

series for xe

( x +1)

xn
, which of the following is a power
n!

( x + 1) 2 ( x + 1) 3
+
+L
a.) 1 + ( x + 1) +
2!
3!
x( x + 1) 2 x( x + 1) 3
+
+L
b.) x + x( x + 1) +
2!
3!
( x + 1) 3 ( x + 1) 4
c.) ( x + 1) + ( x + 1) 2 +
+
+L
2!
3!
x4 x5
d.) 1 + x 3 +
+
+L
2! 3!
x( x + 1) 3 x( x + 1) 4
2
e.) x + x( x + 1) +
+
+L
2!
3!

139

________________________________________________________________________
Questions 16 and 17 refer to the graph below:

16.) The above diagram is a graph of the function f (x) on the interval [-2, 10]. If
g ( x) =

f (t )dt , how many roots does g(x) have on this interval?

0

a.) 4, b.) 3, c.) 2, d.) 1, e.) 0

________________________________________________________________________
17.) Where on the interval [-2, 10] is g (x) increasing?
I.) [-2, 0], II.) [0, 3], III.) [3, 5], IV.) [5, 10]
a.) I only, b.) II only, c.) I and III, d.) II and IV, e.) IV only

________________________________________________________________________
18.) A spherical balloon is inflated with helium at a rate of 200 ft3/min. How fast is the
balloons radius increasing at the instant the radius is 8ft?
a.)

32 ft
25 ft
, b.)
,
25 min
32 min

c.)

75 ft
25 ft
20 ft
, d.)
, e.)
32 min
256 min
25 min

140

________________________________________________________________________

19.) Consider the infinite geometric series

2k . For what value of k will this series

n

n =0

converge to a value of 3?
a.)

3
2
3
, b.) , c.) , d.) The series diverges, e.) Not enough information provided.
2
3
4

________________________________________________________________________
20.) The following graph represents the acceleration of an object:

Which of the following graphs could represent the position of the object?

141

a.)

b.)

c.)

d.)

e.)

21.) What is the average value of the function f ( x) = x 3 sin x on the interval [0, ] ?
a.) 3 6 1 , b.) 3 + 6 , c.) 2 6, d.) 3 6 , e.) 3 + 6

________________________________________________________________________
142

22.) The length of a curve on the interval [0, 5] is given as

1 + 20 x 8 dx. If this

particular curve contains the point (1, 3), which of the following could be an equation for
this curve?
2 5 5 14
x +
5
5
20 9 7
y=
x +
9
9
2 5 5 15 2 5
y=
x +
5
5
2 5 5 15
y=
x +
5
2 5
4
y = 2x + 1

a.) y =
b.)
c.)
d.)
e.)

________________________________________________________________________
23.) A rectangle is to be inscribed between the curve y = 4 2 x 2 and the x-axis, with its
base on the x-axis. For which value of x will the area of this rectangle be maximized?
a.)

3
, b.)
2

2
, c.)
3

3
, d.)
2

2 , e.)

2
2

________________________________________________________________________

143

24.) Calculate the area between the curves y = e x and y = e 3 x , bounded by the lines y =
1 and y = 4.
8
10
a.) ln 4
3
3
8
b.) ln 4 4
3
8
c.) ln 4 2
3
1
2
d.) e12 e 4 + e
3
3
1 12
2
e.) e e 4 e
3
3

________________________________________________________________________
25.) What is the coefficient of x 3 in the Maclaurin series for
a.)

e2x
?
2

2
1
1
1
3
, b.) , c.) , d.) , e.)
3
6
2
3
6

________________________________________________________________________

144

26.) Which of the following represents the derivative of 2 xy = x + y in polar form?

a.) r cos + r sin
1 2r sin
b.)
2r cos 1
c.) 2r cos + 2r sin 1
r cos 2
d.)
sin 1
e.) cos sin 2

________________________________________________________________________
27.) Consider the region bounded by the curve y = x 2 + 3, the line x = 1 , the x-axis, and
the y-axis. What is the volume of the solid generated when this region is revolved about
the y-axis?
a.) 88 , b.) 176 , c.)

24
3
17
, d.) , e.)
5
7
5

________________________________________________________________________

145

28.) What is the radius of convergence of the power series

n =0

a.) (k 3) < x < (k + 2)

b.) ( k 3) > x > (k 3)
c.) (k 3) > x > (k 3)
d.) < x <
e.) Divergent

k n +1
, where k > 0?
( x + 3) n

________________________________________________________________________
End of Section I, Part A
If you finish before time is called, you may check your work in this section ONLY.
Do not proceed to Section I, Part B until told to do so.

146

Section 1, Part B
50 minutes
17 questions
Directions: Solve each of the problems in this section, using the available space for
scratchwork. Select the letter that corresponds to the best solution to the problem. No
credit will be awarded for anything written in the test booklet. Use your time wisely. A
graphing calculator is required for this section of the exam.
Note: 1.) Some of the numerical answers will not be exact. When this is the case, choose
the answer that best approximates the exact numerical value. 2.) Unless otherwise
indicated, the domain of a function f is the set of all real numbers x for which f (x) is a
real number.
29.)

x3 x dx =
2

3
ln 3
3x
ln 3 x 2
3x
+ C , c.)
3 , d.)
+ C , e.) 2 x3 x x + C
a.) 3 + C , b.)
2
2
2 ln 3
x

________________________________________________________________________
t5 +1
gallons per minute.
2t + 3t 2 + 2
How much water has passed through the turbine after one hour?
30.) The flow of water through a turbine is modeled as

a.) 6.4800 10 7 gallons

b.) 1.0639 10 6 gallons
c.) 7.7762 10 9 gallons
d.) 7.7761 10 8 gallons
e.) 6.5001 10 6 gallons
________________________________________________________________________

147

31.) Given the infinite sequence u (n) =

a.) 0, b.) , c.) ln 3 , d.)

en + 1
, what is lim u (n) ?
n
3n 1

1
, e.) ln(3 + e)
ln 3

________________________________________________________________________
32.) Consider the graph of the function f (x) below:

Which of the following must be true?

I.) f ( x) < 0 on (a, b) , II.) f (a) does not exist, III.)

f ( x)dx exists
b

a.) I only, b.) III only, c.) I and II, d.) I and III, e.) I, II, and III

________________________________________________________________________
33.) Which of the following represents the value of the approximation y(1.2) of the curve
y = x ln x using a 4th-order Taylor polynomial centered at x = 1?
1631
409
1661
368
219
a.)
, b.)
, c.)
, d.)
, e.)
7500
1875
7500
1682
1000

________________________________________________________________________
148

dF
= 1.77t 2 + e 0.03t models how quickly the rate of a
dt
quantity changes. If the initial rate at which the quantity changes is 0, how much of the
quantity is there when t = 10 ?
34.) The differential equation

a.) 599, b.) 1187, c.) 589, d.) 1520, e.) 1177

________________________________________________________________________
35.) The approximate value of the area under the curve y = 2 x 2 2 x 3 + 2 x + 6 on the
interval [0, 1] can be found by using the trapezoidal rule. When using 5 trapezoids, what
is the approximate area under the curve on this interval using only the trapezoidal rule?
a.) 7.17, b.) 7.167, c.) 7.16, d.) 7.2, e.) 7.1667

________________________________________________________________________

5 x3

36.) Given f ( x) = sin tdt , f (x) =

1

a.) 15 x sin(5 x )
b.) 15 x 2 cos(5 x 3 )
c.) cos(5 x 3 ) + cos(1)
d.) 5 x 3 sin(15 x 2 )
e.) 5 x 3 cos(15 x 2 )

149

________________________________________________________________________
d2y
37.) If the curve y(x) is defined parametrically as x = 3t 3 , y = 5 cos t ,
=
dx 2
a.) 5 sin t
45t 2 cos t + 90t cos t
b.)
729t 6
45t 2 cos t + 90t
c.)
81t 2
5 sin 2 t
d.)
9t 2
sin t cos t (9t 3 + 18t 2 )
e.)
9t 6

________________________________________________________________________
38.) Determine the area inside the circle r = 1 and outside the cardioid r = 3(1 + sin ) .
a.) 1.334, b.) 12.412, c.) 6.206, d.) 8.315, e.) 0.667

________________________________________________________________________
39.) Which of the following are -values at which there are horizontal tangents of
r = 4 sin on the interval [0, 2 ] ? (Note: Not all of the values may be listed.)
I.) = 0, II.) = , III.) 5.6642
a.) I only, b.) II only, c.) III only, d.) I, II, and III, e.) None of the above

150

________________________________________________________________________
40.) On the interval [0, 5], which of the following graphs could represent the function f
1 5
with the property that
f (t )dt = 2 ?
5 0
a.)
b.)

c.)

d.)

e.)

________________________________________________________________________
41.) Sand is poured on top of a conical pile at the rate of 200 ft 3 / sec , causing the radius
1
of the base of the pile to increase at ft/sec. How quickly is the height of the pile
2
increasing at the instant the radius of the base is 18 ft and the height of the pile is 10 ft?

151

a.) 0.2778 ft/sec

b.) 57.2958 ft/sec
c.) 0.0339 ft/sec
d.) 0.0617 ft/sec
e.) 16.2001 ft/sec

________________________________________________________________________
42.) The rate of change of a quantity is given by the differential equation
dy
6t + 3
= 2
. If y (2) = 7, which of the following is the value of the constant term in
dt t + 2t 3
the equation for y?
a.) 7

15
9
15
9
15
ln 5 , b.) 7 ln 5 , c.) 7 ln 2 , d.) 7 ln 2 , e.) 7 ln 7
4
4
4
4
4

________________________________________________________________________
43.) If the velocity of a particle is given by the parametric equations
x(t ) = 3t 3 , y (t ) = sin 2 t , what is the speed of the particle at t = 1?
a.) 3.148, b.) 9.909, c.) 3.708, d.) 1.926, e.) 9.046

________________________________________________________________________

152

44.) Consider the graph of the function f below:

Which of the following must be true?

I.)

f ( x)dx does not exist,

0

II.) lim+ f ( x) = lim f ( x), III.) f (A) does not exist

x A

x A

a.) I only, b.) II only, c.) III only, d.) I and II, e.) II and III

________________________________________________________________________
45.) If y = 5 x 66 ,

d 66 y
=
dx 66

a.) 0, b.) 5(66!), c.) 330(66!), d.) 5(66!)x, e.) 330(64!)

________________________________________________________________________
End of Section I, Part B
If you finish before time is called, you may check your work in this section ONLY.
Do not proceed to Section II, Part A until told to do so.

153

Section II, Part A

45 minutes
3 problems
Directions: Solve each of the problems in this section, showing all work leading to your
answer. Credit will be awarded based upon the accuracy and completeness of your work
and answers. Unless otherwise indicated, algebraic or numerical answers need not be
simplified. If decimal approximations are used, be sure to give the final answer to three
decimal places after the decimal point. Use your time wisely. A graphing calculator is
required for this section of the exam.
Note: Unless otherwise indicated, the domain of a function f is the set of all real
numbers x for which f (x) is a real number.

1.) Let the shaded region shown above be enclosed by the graphs of y = sin( x 2 ) and
y = x3 2x 2 .
a.) Calculate the area of the shaded region.

b.) Determine the volume of the solid generated when the shaded region is revolved
about the y-axis.

154

c.) Let the shaded region be the base of a solid. For this solid, each cross section
perpendicular to the x-axis is a semicircle. Calculate the volume of this solid.

________________________________________________________________________
2.) Two particle are moving in the xy plane. The velocity of Particle 1 is given
parametrically as:
v x = t 3 5t 2 3t
1

vy = e 2 .
The velocity of Particle II is given as v x = t 2 3, while v y is not explicitly given.
Both particles are initially at the origin.
a.) Determine x(t ) and y (t ) for Particle 1 and x(t ) for Particle 2.

155

b.) It has been found that y (t ) is an increasing linear function for Particle 2. If Particle 1
and Particle 2 have the same y-coordinate at t = 0.708399, find y (t ) for Particle 2.

c.) Find the acceleration vector of each particle when t = 1.5.

d.) Do any of the particles change direction from t = 0 to t = 2 ? If so, find the total
distance traveled by the particle(s) when it (they) first change direction.

156

________________________________________________________________________
3.) Consider the differential equation

dy
= x + 2xy 2 .
dx

a.) Using Eulers method, approximate the value of y (1.4) by using increments of 0.1 if
y (1) = 0.5 .

b.) In the grid provided, draw the slope field that represents the family of solutions for
this differential equation.

157

y
3

x
3
c.) Rewrite in polar form the expression to which

dy
is equal.
dx

________________________________________________________________________
End of Section II, Part A
If you finish before time is called, you may check your work in this section ONLY.

158

Do not proceed to Section II, Part B until told to do so.

Section II, Part B
45 minutes
3 problems
Directions: Solve each of the problems in this section, showing all work leading to your
answer. Credit will be awarded based upon the accuracy and completeness of your work
and answers. Unless otherwise indicated, algebraic or numerical answers need not be
simplified. Use your time wisely. No calculator of any kind may be used in this section.
Note: Unless otherwise indicated, the domain of a function f is the set of all real
numbers x for which f (x) is a real number.

4.) The graph above represents the rate at which the altitude of a traveling object changes.
Time t is on the horizontal axis and change in altitude A/t is on the vertical axis. The
graph is composed of line segments and a semicircle.
a.) On what intervals, if any, is the object gaining altitude? Justify your answer.

b.) What is the height of the object when t = 5?

159

c.) On a set of axes (not provided), sketch the graph of the objects acceleration (A/t2)
versus time. The graph does not have to be to scale. Justify your sketch.

d.) Calculate the average rate of change of altitude of the object from t = 0 to t = 10.

160

________________________________________________________________________
5.) The following are true of the functions of f ( x) and g ( x) :
i.) f ( x) = x 2 g ( x)
f (1 + k ) f (1)
ii.) lim
=1
k 0
k
iii.) lim f ( x) = 2
x

iv.) f (1) = 0, f (1) = 1

a.) Does the improper integral

( x 2 g ( x))dx converge or diverge? If it

converges, to what value does it converge?

b.) Show that g (1) = 0.

c.) Using LHpitals Rule, show that lim

x 1

f ( x)
= 1.
g ( x)

161

6.) A function f, which has a range of [-2, 1], is defined by the following power series for
all real numbers x:

n =0

(1) n 2 x n n
2x 4x 2 6x 3
=
+

+L.
(3n + 1)!
4!
7!
10!

a.) Find f (0) and f (0) . Does f (x) have a minimum, maximum, or point of inflection
at x = 0? Justify your answer.

b.) What minimum order polynomial expansion centered at x = 0 must be used to

1
approximate f (1) with an error no greater than
?
100

162

c.) Write the 4th-order Maclaurin polynomial for f ( x 3 + 2) .

________________________________________________________________________
End of exam
If you finish before time is called, you may check your work in this section and Section II
A, but you will not be allowed the use of a calculator.

163

Explained Answers to the Practice Exam

Section I A:
1.) C. This problem involves implicit differentiation:
3 xy = x 4 y 3
dy
dy

3 x + y = 3 x 4 y 2
+ 4x3 y 3
dx
dx

dy
dy
3x + 3 y = 3x 4 y 2
+ 4x3 y 3
dx
dx
dy
4 2
(3 x 3 x y ) = 4 x 3 y 3 3 y
dx
dy 4 x 3 y 3 3 y
.
=
dx 3 x 3 x 4 y 2
2.) B. The carrying capacity is reached when the population stops growing. Since the
equation given represents the rate of change, one sets it equal to zero and solves for N.
N

The whole equation will equal zero when the term 1 equals zero.
72
3.) B. This is a definite integral in which one of the limits of integration is a variable:

3 sin xdx = 3 cos x = ( 3 cos(t ) ) ( 3 cos( ) ) = 3 cos t 3.

t

4.) B. For a volume of solid of revolution problem, one has three possibilities: the disc
method, the washer method, or the shell method. The disc method would not work for
this problem because the radius would need to be perpendicular to the y-axis and, due to
the way in which the functions intersect, there would be no direct way to express the
radius algebraically. The washer method cannot be used because there is no hole in the
resulting solid. One must use the shell method. The formula associated with this method

when revolution is about the y-axis is V = 2 hr dx , where h is the height of the shell
a

(the distance between the bounding functions) and r is the radius of the shell (the distance
from the axis of rotation to the bounded functions; for a rotation about the y-axis, this
distance is simply x). Thus, the volume should be expressed as
2

5.5

(g ( x) f ( x) )

xdx .

5.) A. The question is asking for the slope of the tangent line, which is

dy
. For a set of
dx

dy
dy
3t 2
3
dy dt
=
= t 2 sec t.
parametric equations,
. Thus, in this instance,
=
dx 5 cos t 5
dx dx
dt

164

6.) E. This integral requires integration by partial fractions. Factoring the denominator,
one of the factors is linear, while the other is cubic, meaning that the numerators of the
partial fractions will be constant and quadratic, respectively:
Bx 2 + C
2x 3 + 2x 2 + 1 A
= + 3
x x + 2x 2 + 1
x 4 + 2x 3 + x
To solve for A, B, and C, one multiplies each side of the equation by the denominator:
2 x 3 + 2 x 2 + 1 = A( x 3 + 2 x 2 + 1) + Bx 3 + Cx
2 x 3 + 2 x 2 + 1 = Ax 3 + 2 Ax 2 + A + Bx 3 + Cx
2 x 3 + 2 x 2 + 1 = x 3 ( A + B) + x 2 (2 A) + x(C ) + A
Since the constant on the left is 1, this must also be the value of A, since it is the only
constant on the right. Since there are no terms with x on the left, this must also be the
case with the right, so C = 0. Finally, since the coefficient of x 3 on the left is 2, this must
also be the case on the right, meaning A+B = 2. Since A was already found to be 1, B
must also be 1. One can now integrate:

2x3 + 2x 2 + 1
dx =
x 4 + 2x3 + x

= ln x + ln x +

dx
+
x

x2
dx =
x3 + 2x 2 + 1

dx
+
x

x2
dx +
x3

x2
dx +
2x 2

x dx
2

1
1
1
1
x + x 3 + C = 2 ln x + x 3 + x + C.
2
3
3
2

7.) C. The best way to do this problem is to test each series separately. The first series
diverges by the nth term divergence test. Since the degree in the numerator and
denominator is the same, and since the limit is being taken at infinity, one takes the ratio
of the leading coefficients. Since this ratio is 1 in this case, one can conclude that the
series diverges because the limit is a number other than 0. The second series diverges by
2(n + 1)! 10 n
2(n + 1)n! 10 n
n +1
lim
the ratio test: lim
=
= lim
= . Since this limit is
n
n
+
+
(
1
)
(
1
)
n 10
2n! n 10
2n! n 10
greater than 1, the series diverges. The third series converges by the comparison test.

3
. Since r < 1, the geometric
4
n =1
series converges. Since this larger series converges, the original smaller series must also
converge. Thus, the third is the only convergent series.
One can compare it to the larger geometric series

8.) C. This problem just requires differentiation with the chain rule:
2

1
1
1
d 3
ln(3x) + 5 = [ln(3x) + 5] 3 =
.
3
dx
x 3 x 3 [ln(3x) + 5] 2

165

9.) C. This problem requires integration by parts. Since one term is a polynomial and the
other is an easily integrable function, one can use the tabular method:
u
3x 3
9x 2
18x
18
0

1
+1
-1
+1
-1
+1
-1

dv
cosx
sinx
-cosx
-sinx
cosx

From this, one finds that the integral is 3x 3 sin x + 9 x 2 cos x 18 x sin x 18 cos x + C .
10.) D. In order for the function to be continuous, the second part must begin where the
first part ended. The y-value at which the first part ended may be determined by
substituting 2 for x: (2) 2 1 = 3. Thus, the second part, which is a horizontal line, must
5
have a value of 3. Thus, one can find a: 3a 2 = 3 a = .
3
11.) E. While the graph of the derivative crosses the x-axis 6 times and is non-existant at
one point, this only means that the actual function has 7 critical points on the interval.
These critical points indicate relative extrema only. The function could also have
absolute extrema at the endpoints. Thus, a conclusion cannot be reached without more
information on the function.
12.) A. Attempting to substitute infinity for initially results in the indeterminate form
1
1
ln 2

1
1

1
0 0. Invoking a property of logarithms, lim 2 lim ln 2 = lim = .

1
ln 2
2

Thus, one uses LHpitals Rule: lim = lim = 0. This is not the answer,

however. Recall that the natural logarithm of the limit was taken. Thus, one must
reverse this by exponentiating: e 0 = 1. This is the final answer.
13.) C. In order to find the value that satifies the Mean Value Theorem, one must set the
slope of the tangent line (the derivative) equal to the slope of the secant line (the average
slope over the interval). For this problem,
y y (2) y (0) 13 1
msec =
=
=
= 6.
x
20
2
One then sets this value equal to the derivative to find at what value of x the derivative
equals this average slope:

166

dy
4
2
= 3x 2 + 2 = 6 x 2 = x =
. This is the value of c that satisfies the Mean
dx
3
3
Value Theorem.
14.) A. This integral can be manipulated to be in the form

du
, which equals
+ a2

1
u
tan 1 . This can be achieved by multiplying the integrand by 3 and multiplying the
a
a
1
outside by :
3
1
3

3dx
4 + 9x 2

= lim

a 1

11
3x
3x
3a 1
3
1
1
= lim tan 1 = lim tan 1 = lim tan 1
tan 1
2
a 3 2
a 6
2
2 1 a 6
2 6
2
4 + 9x
1
3dx

Note that the inverse tangent of infinity is

(because tan x is undefined at that value).

3
3 1
tan 1 .
=
2
2 12 6

3a 1
1
Thus, lim tan 1
tan 1
a 6
2
6

15.) B. In this problem, one is able to substitute x + 1 wherever there is an x in the series
and then multiply every term by x. Note that this exact procedure must be taken; one
cannot multiply every term by x and then substitute x + 1 for every x (This would actually
yield the answer represented by choice C). The reason for this is because one must
consider the manipulation of the functions argument first and only then consider what is
being multiplied. Thus, one must first consider the x in e x by replacing every x with
x + 1, and then multiply every term by x.
16.) E. The function g ( x) =

f (t)dt represents an accumulation function. By calculating

0

the area (including the sign negative or positive) under f (x), one can determine the
values of g (x). The areas associated with the graph are shown below:

Note that the left-most triangle has an associated area of -4 because the accumulation
function begins at x = 0. Thus, one must negate whatever area is to the left of that.
Starting at x = -2 and summing all of the areas up to x = 10, one would see that a value of
zero is never reached: 4 4 + / 2 / 2 + 6 = 2 . Thus, g (x) has no roots on the
interval.
167

17.) C. Differentiating both sides of g ( x) =

f (t )dt , one can see that f (x) is the

0

derivative of g (x). Thus, the graph depicted represents the derivative of g (x). Since this
is the case, f (x) would be increasing on whichever interval g (x) is positive. The only
intervals on which f (x) is always positive are [-2, 0] and [3, 5]. Thus, f (x) is increasing
on these intervals.
18.) B. This is a related rates problem. Perhaps the best way to go about solving it is to
record all of the information given in the problem, write mathematically what the
problem is asking for, and write the equation to be used. In this problem, it is given that
dr
dV
ft 3
= 200
. One must find
when r = 8ft. The equation to be used is the
min
dt
dt
4
expression for the volume of a sphere: V = r 3 . Unfortunately, it is not guaranteed that
3
these sorts of geometric formulae will be given on the exam. One now differentiates the
equation, substitutes the known values, and solves for the unknown element:
d
d 4
dV
dr
dr

(V ) = r 3
= 4 r 2
(200 ft 3 / min) = 4 (8ft ) 2
dt
dt 3
dt
dt
dt

dr
256 ft 2
200ft 3 / min
dt dr = 25 ft .
=
2
2
dt 32 min
256 ft
256 ft

19.) A. This geometric series can be rewritten in the more familiar form

n =0

1
k .
2

1
< 1, the series certainly converges. This is a geometric series, so its sum has
2
a
. Since S n and r are known, one can find a:
the formula S n =
1 r
a
3
(3) =
a= .
2
1
1
2

Since r =

20.) D. This is a very concept-based question. It concerns the relationship between the
graph of a function (position) and its second derivative (acceleration). For this kinds of
problems, it is best to observe each part of the graph separately and conclude such
relationships. Beginning with the first part of the acceleration graph, the acceleration
linearly decreases until it reaches zero. On the velocity graph, one would see a quadratic
increase with linearly decreasing slope until a relative maximum is reached. On the
position graph, one would see a cubic increase until a point of inflection is reached.
Thus, one can rule out choices E and B because they do not exhibit a point of inflection at
the value where the acceleration is zero. Analyzing the acceleration graph further, the

168

acceleration is negative and linearly decreasing until it becomes constant. After this
constant value, the acceleration begins to increase toward zero. One the velocity graph,
one would see a quadratic decrease in velocity with linearly decreasing slope and then a
linear decrease in velocity, followed by a quadratic decrease in velocity with linearly
increasing slope. On the position graph, one would see a cubic increase in position with
quadratically decreasing slope until another point of inflection is reached. This leaves
only choice D.
21.) C. Recall that the average value of a function over a certain interval [a, b] has the
b

f ( x)dx

following formula: average value =

. The integral in this question requires
a

ba
integration by parts. However, since one term is a polynomial and the other is an easily
integrable function, one can use the tabular method:
u
x3
3x 2
6x
6
0

dv
sinx
-cosx
-sinx
cosx
sinx

1
+1
-1
+1
-1
+1
-1

From this,

x sin xdx = x cos x + 3x sin x + 6x cos x 6 sin x + C = cos x(6x x ) + sin x(3x
3

6) + C.

One can know perform the appropriate definite intregral:

x 3 sin xdx

[cos( )(6

) + 0 [cos(0)]

= 2 6.

22.) C. Recall that the arc length of a curve has the formula l =

b
a

dy
1 + dx. With
dx

dy
this formula in mind, one can see that = 20 x 8 for this problem. Thus, this
dx
problem is essentially asking one to solve a differential equation and use a set of initial
conditions:
dy
2 5 5
= 20 x 4 = 2 5 x 4 y = 2 5 x 4 dx =
x + C.
dx
5

169

2 5 5
15 2 5
(1) + C = 3 C =
.
5
5
2 5 5 15 2 5
.
Thus, y =
x +
5
5
It is given that y (1) = 3, so

23.) B. This is an optimization problem, which involves setting the derivative of a certain
equation equal to zero and finding the values that make the original function the largest
or smallest that it can be on a certain interval. In this problem, one wishes to find the
value that makes the area function have the greatest magnitude. To do this, one
determines the derivative of the area function, sets this expression equal to zero, and finds
the value at which the derivative shifts from being positive to being negative (i.e. where
the area function has a relative maximum). The area in question is that of a rectangle
with its height equal to the distance from the x-axis to the function, which is merely
y = 4 2 x 2 , and its length equal to 2x, since the base is symmetric about the y-axis.
Thus, A = lh = (2 x)(4 2 x 2 ) = 8 x 4 x 3

dA
2
= 8 12 x 2 = 0 x = . Since at
dx
3

2
the derivative shifts from positive to negative values, the area function has a
3
relative maximum at this point.
x=+

24.) C. Since no calculator may be used on this portion of the exam, one must use a little
mathematical intuition here. Since y = e 3 x increases more quickly than y = e x , the
former curve will generally lie closer to the y-axis than the latter curve. Thus, y = e 3 x is
farther to the left and y = e x is farther to the right. Since the boundaries in this case are
horizontal lines, there is no top and bottom function that one can use to determine the
area. However, there is a left and right curve, but one must have everything,
including the limits of integration, in terms of y since the infinitesimal rectangles are now
perpendicular to the y-axis instead of the x-axis. The leftmost curve in terms of y is
ln y
x=
and the rightmost curve in terms of y is x = ln y. One now evaluates the definite
3
4
4
4
ln y
1 4

integral: A =
" Right"" Left" dy =
ln
y

dy
=
ln
y
dy

ln y dy. To

3
3 1
1
1
1
integrate the natural logarithmic function, one technically needs to use integration by
parts. However, I recommend memorizing what the integral is to save time:
8
1
4
4
ln y dy = y ln y y + C . Thus, A = ( y ln y y ) 1 ( y ln y ) 1 = ln 4 2.
3
3

25.) A. This question concerns the algebraic manipulation of power series. Note that the
Maclaurin series for e x is one of the power series that should be memorized for the exam.
The others can be found on page 61. For this problem, one must substitute 2 x for x and

170

1
multiply by . Since the Maclaurin series for e x is
2
the Maclaurin series for

1 2x
e is
2

n =0

n =0

xn
x2 x3 x4
= 1+ x +
+
+
+ L,
n!
2! 3! 4!

(2 x)
1
4x 2
8 x 3 16 x 4
= +x+
+
+
+ L. Thus, the
2n!
2
2 2! 2 3! 2 4!
n

8
2
coefficient of x 3 is
= .
2 3! 3
26.) B. One must first differentiate this expression implicitly:

d
d
dy
dy
dy

dy
(2 xy) =
( x + y ) 2 x + y = 1 +
2x + 2 y = 1 +
dx
dx
dx
dx
dx

dx
dy
dy 1 2 y
(2 x 1) = 1 2 y
.
=
dx
dx 2 x 1
dy
in polar form, one must recall that x = r cos and y = r sin . So,
dx
dy 1 2 y 1 2(r sin ) 1 2r sin
=
=
=
.
dx 2 x 1 2(r cos ) 1 2r cos 1

To express

27.) E. The function in question is a parabola that intersects the y-axis at y = 3. It is

bounded between x = 0 and x = 1. One must decide whether to use the disc method, the
washer method, or the shell method. The disc method cannot be used because the area is
being revolved about the y-axis and there is no direct way to express the radius of a disc
in terms of the function given (If it were revolved about the x-axis there would be, since
the radius would merely equal the distance from the x-axis to the function). The washer
method need not be used because there is no hole in the solid of revolution. Thus, one
uses the shell method. The formula associated with this method when revolution is about

the y-axis is V = 2 hr dx , where h is the height if the shell (the distance between the xa

axis and the function) and r is the radius of the shell (the distance from the axis of
rotation to the bounded function; for a rotation about the y-axis, this distance is simply x).
Thus, the formula for the volume of this solid is
1

17
3
1
.
V = 2 ( x + 3) xdx = 2 x 5 + x 2 =
5
2 0
5
0

28.) B. Determining the radius of convergence of a power series requires using the ratio
test and determining for which values of x the series converges. Perhaps the only more
slightly complicated aspect of this problem is that the constant k is unknown.

171

n =0

k n+ 2

n +1
k n +1
( x + 3)
lim
n
( x + 3) n
k n +1

n
( x + 3)

k n + 2 ( x + 3) n

= lim
n ( x + 3) n +1
k n +1

k
=
.
x+3

a n +1
< 1. Thus,
n a
n

The ratio test states that a series will converge if lim

k
k
x+3
< 1 1 <
< 1 1 >
>1
x+3
x+3
k

k > x + 3 > k ( k 3) > x > (k 3). This represents the radius of convergence.

Section I B:
29.) D. The integral of a general exponential a u is given as

a du = ln a a
u

+C.

For x3 x dx, one must first multiply the integrand by 2 to give du = 2 xdx and
multiply the outside of the integral by

1
to compensate for this. The solution to the
2

1 1 x2 3 x
3 =
.
integral is, thus:
2 ln 3
2 ln 3
30.) B. Recall that the definite integral of a rate yields a total amount. Since the equation
in this problem represents flow, integration over a certain time period will yield the total
amount associated with that time period. One should perform the following definite
60
t5 +1
integral on the calculator:
dt = 1063878.5 gallons.
2
0 2t + 3t + 2

31.) A. No tests are really required for this problem. In both the numerator and
denominator of this fraction are exponentials. However, the denominator contains an
exponential whose base (3) is larger than that in the numerator ( e 2.71828 ). Thus, the
denominator will increase more quickly than the numerator; the addition or subtraction of
1 is negligible. The effect of the more rapid increase in the denominator is that the
fraction approaches zero as n becomes larger and larger.
32.) C. Only the first two statements are correct. The first statement is correct because
on the open interval (a, b), the function is decreasing, meaning the derivative is negative
on that interval. The second statement is correct because at a, the function exhibits a
corner, meaning the function is not differentiable at that point. The third statement is not
necessarily true. While it does seem that the graph will continue to decrease as x
becomes larger and larger, implying that the area under the curve converges on a finite

172

value, this is not guaranteed. Without more knowledge of the functions behavior, one
cannot truly know whether the function continues to decrease toward a y-value of 0.
33.) C. One must find the first several terms of the polynomial expansion for y = x ln x
to approximate y (1.2). Recall that the formula for a Taylor expansion is

n =0

f ( n ) (c)( x c) n
. It is given that c = 1. Applying the formula to y = x ln x up to n = 4,
n!

(1 /(c))( x c) 2 (1 /(c) 2 )( x c) 3 (2 /(c) 3 )( x c) 4

x ln x [(c) ln(c)] + [(1 + ln(c))( x c)] +
+
+

2!
3!
4!

Substituting 1 for c and 1.2 for x,

(1 /(1))(0.2) 2 (1 /(1) 2 )(0.2) 3 (2 /(1) 3 )(0.2) 4
(1.2) ln(1.2) [(1) ln(1)] + [(1 + ln(1))(0.2)] +
+
+

2!
3!
4!

1661
0.2 + 0.02 0.001333 + 0.0001333
.
7500

34.) D. This question requires careful reading. It states that the equation given models
how quickly the rate of a quantity changes. Thus, it is actually measures the rate of rate!
This means that the equation given represents the second derivative of the equation
representing the amount, meaning one must integrate twice. Integrating the first time
yields the equation representing the rate of change of the quantity, or the first derivative
of the equation representing the amount. It is stated that the initial rate is zero. With
these initial conditions, one may find the exact equation representing the first derivative:
First derivative =

(1.77t

+ e 0.03t dt = 0.59t 3 33.333e 0.03t + C

0.59(0) 3 33.333e 0.03( 0 ) + C = 0 C = 33.333.

To determine the amount at t = 10, one integrates the equation for the first derivative
from t = 0 to t = 10. Using the graphing calculator:

10

(0.59t 3 33.333e 0.03t + 33.333)dt = 1520.3536.

35.) C. Recall the general formula for the trapezoidal rule:

1
A = w( f (a ) + 2 f ( x1 ) + 2 f ( x 2 ) + L + 2 f ( x n ) + f (b) ) , where w is the width of each
2
trapezoid (found by dividing the length of the interval by the number of required
trapezoids) and a and b are the endpoints of the interval. For this problem,

173

A =

1
2

(0.2)( f (0) + 2 f (0.2) + 2 f (0.4) + 2 f (0.6) + 2 f (0.8) + f (1)) =

(0.2)(6 + 12.928 + 13.984 + 14.976 + 15.712 + 8) = 7.16

2
36.) A. While one could execute the steps of evaluating the definite integral, it is helpful
to use a property of accumulation functions. When taking the derivative of an
accumulation function with a constant lower limit of integration and a variable upper
limit of integration, one simply substitutes the upper limit for every t and multiplies this
expression by the derivative of the upper limit. Thus,
3

d 5x
f ( x) =
sin tdt = 15 x 2 sin(5 x 3 ) .
dx 1

37.) B. Recall that the formula for the second derivative of a parametrized curve is:
d dy

2
dy
d y dt dx
=
. One first determines
:
2
dx
dx
dx
dt
dy
dy dt 5 sin t
. One then differentiates this expression with respect to t:
=
=
dx dx
9t 2
dt
d dy (9t 2 )(5 cos t ) (5 sin t )(18t ) 90t sin t 45t 2 cos t
=
=
dt dx
81t 4
81t 4

One finally divides this expression by

dx
:
dt

d 2 y 90t sin t 45t 2 cos t 1

90t sin t 45t 2 cos t
=
=
dx 2
81t 4
9t 2
729t 6
38.) E. This problem should be done via graphing calculator. The two polar curves are
shown below and the shaded region represents the area to be found:

174

To determine where the curves intersect, one sets them equal and solves for :
2
= 0.72972766. This angle is measured in a
3
clockwise sense. To express it in a counterclockwise sense, one adds it to 2 (i.e.360 o ) .
Thus, one angle is 2 0.72972766 = 5.553458. Since the value of sine is negative in
quadrants III and IV, the other angle must be in quadrant III, 0.72972766 radians from
(180 o ) . Thus, the second angle is + 0.72972766 = 3.871320. Now that the two
limits of integration have been found, one can compute the area:
1 = 3 + 3 sin sin =

1 2
1 5.553458
1
r2 r12 d =
(3 + 3 sin ) 2 1 d = (1.333963) 0.6669815
2
2 3.871320
2
(Note that though the integral is negative, area must always be positive).
A=

dy
dx
= 0 and where
0 at the
d
d
same time. One first determines the derivative of y with respect to :

39.) B. Recall that a horizontal tangent exists wherever

y = r sin = 4 sin (sin ) = 4 sin 2

dy
= (4 )(2 sin cos ) + (sin 2 )(4) = 8 sin cos + 4 sin 2
d

The zeros of this equation can be found by graphing it on the calculator in

FUNCTION mode and determining where it crosses the x-axis. Doing this, one finds
dx
does
zeros at: = 0, 1.8365972, , 4.8158423, 2 . One must now make certain that
d
not also have zeros at these -values:
x = r cos = (4 sin ) cos = 4 sin cos (Let 4 sin be the " first term" )
dx
= (4 sin )( sin ) + (cos )(4 cos + 4 sin ) (Product rule within product rule)
d
= 4 sin cos + 4 cos 2 4 sin 2 .

The zeros of this function are: = 0, 0.72024122, 2.8681853, 3.4159489, 6.0864406 .

dy
dx
Since both
and
are zero when = 0, one can make no conclusion about this
d
d
dy
is zero, however, do have horizontal tangents.
point. The other -values where
d
However, the only one specified in the choices is , making this the only correct choice.
40.) E. The expression given is the average value function. It states that on the interval
[0, 5], the average value of the function is 2. To decide which graph has an average value
of 2 on this interval, one finds the total area under each curve (including any negative
175

signs, since the average value concerns the integral) and divides this value by 5. The
only graph whose total area divided by 5 is 2 is the one depicted in choice E.
41.) C. This is a related rates problem. As always, one should record every quantity
given, determine what must be found, and find the equation to be used. In this question,
dV
dr 1
= 200ft 3 / sec,
= ft/sec, and that when r = 18ft, h = 10ft. The
it is given that
dt
dt 2
dh
at these values. For this problem, one can use the
question is asking for the value of
dt
1
equation for the volume of a cone: V = r 2 h . Do not expect this formula to be given
3
on the exam; it is better to memorize it and those for other geometric solids. One
differentiates this equation, substitutes all necessary values, and solves for the unknown:

d
d 1
dV 1 2 dh
dr

(V ) = r 2 h
= (r ) + (h) 2r Note that the product rule was used here.
dt
dt 3
dt 3 dt
dt

1
2 dh
(200ft 3 / min) = (18ft ) + (10ft)(2(18ft)(1/2ft/min) )
3
dt

3
2 dh
3
(200ft / min) = 324ft
+ 180ft / min
3
dt

600 3
2 dh
3
ft / min = 324ft
+ 180ft / min

dt
dh
= 0.0339ft/min.
dt
42.) A. Integration by partial fractions is required to solve this differential equation:
6t + 3
6t + 3
A
B
dt 2
=
+
t + 2t 3
t + 2t 3 (t + 3) (t 1)
6t + 3 = A(t 1) + B(t + 3) = At A + Bt + 3B = t ( A + B) + 3B A
y=

Since the constant on the left side is 3, 3B A = 3 . Since the coefficient of t on the left
side is 6, A + B = 6. Solving these two equations simultaneously, one finds that
15
9
A = and B = . One can now integrate:
4
4
y=

6t + 3
15
dt =
4
t + 2t 3
2

dt
9
+
t +3 4

Given the initial condition y (2) = 7,

B
15
9
= ln t + 3 + ln t 1 + C
t 1 4
4

15
9
15
ln (2) + 3 + ln (2) 1 + C = 7 C = 7 ln 5.
4
4
4

176

dx dy
43.) E. The speed ( v ) of a particle is given parametrically as v = + .
dt dt
Thus, for this particle,
dx
dy
= 9t 2 ,
= 2 sin t cos t.
dt
dt
2
v = 9t 2 + (2 sin t cos t ) 2 = 81t 4 + 4 sin 2 t cos 2 t At t = 1,
v = 81 + 4 sin 2 (1) cos 2 (1) 9.045818.

( )

44.) E. The first statement is not necessarily correct because a definite integral may very
well exist over an interval on which there are discontinuities. The second statement is
true; approaching A from the left side or the right side yields the same limit. The third
statement is also true; since there is a removable discontinuity at A, the derivative does
not exist there. Remember that continuity is a pre-requisite for differentiability.
45.) B. Obviously it would be quite impractical to evaluate 66 derivatives! Instead, one
should evaluate the first several derivatives and find a pattern that can be used to find the
66th derivative:
dy
= 5(66) x 65
dx
d2y
= 5(66 65) x 64
2
dx
d3y
= 5(66 65 64) x 63
3
dx
d4y
= 5(66 65 64 63) x 62
4
dx

Note that the order of the derivative corresponds to how many new terms are
being multiplied. Thus, when the 66th derivative is taken, there will be 66 multiplied
d 66 y
= 5(66!) .
terms. This means that
dx 66
Section II A:
1.) a.)

A=

x2

("Top"-"Bottom")dx
x1

1.8862947

[(sin( x

)) ( x 3 2 x 2 ) dx = 2.1810975 1.181.

177

x2

b.) Shell Method About y - axis : V = 2 rhdx

V = 2

1.8862947

c.) V =

x1

x (sin( x 2 )) ( x 3 2 x 2 ) dx = 2 (2.51117) = 15.77814665 15.778.

x2

(Area of cross section)dx

x1

1 2
1 (sin( x 2 )) ( x 3 2 x 2 )
Semicircle A = r The radius is 1/2 the base A =

2
2
2

1
V =
8

1.8862947

[(sin( x

2
1
)) ( x 3 2 x 2 ) dx = (2.181097499) = 0.856514985 0.857
8

Rubric:
Part A: 1 point for setting up integral correctly
1 point for correct limits of integration
1 point for correct answer to 3 decimal places
Part B: 1 point for specifying shell method
1 point for setting up integral correctly
1 point for correct answer to 3 decimal places
Part C: 1 point for volume equation
1 point for setting up integral correctly
1 point for correct answer to 3 decimal places
Total: 9 points
2.) a.) Particle 1:
dx
dt
1
5
3
3
2
x = (t 5t 3t )dt = t 4 t 3 t 2 + C
4
3
2
1 4 5 3 3 2
1
5
3
(0) (0) (0) + C = 0 C = 0 x(t ) = t 4 t 3 t 2
4
3
2
4
3
2

v x = t 3 5t 2 3t =

1
t
2

dy
dt
1
t
12 t
y = e dt = 2e 2 + C

vy = e

178

2e 2

(0)

+ C = 0 C = 2 y (t ) = 2e 2 2

Particle 2:
dx
dt
1
2
(t 3)dt = t 3 3t + C
3

vx = t 2 3 =
x=

1 3
1
(0) 3(0) + C = 0 C = 0 x(t ) = t 3 3t
3
3
b.) If Particle 2s y-position is an increasing linear function, let y be generally represented
as y (t ) = at + b . Since y (0) = 0, b = 0 .
y-coordinate of Particle 1 at t = 0.708399:
y (0.708399) = 2e

1
( 0.708399 )
2

2 = 0.85007891.

If this is also the y-coordinate of Particle 2 at this time,

6
6
y (0.708399) = a(0.708399) = 0.85007891 a = 1.2 = . Thus, y (t ) Particle 2 = t.
5
5
dv x
= 3t 2 10t 3
dt
dv y 1 12 t
= e
a y (t ) =
2
dt
2
a x (1.5) = 3(1.5) 10(1.5) 3 = 11.25
1
1 2 (1.5)
a y (1.5) = e
= 1.0585
2
a(t ) = 11.25i + 1.0585 j.

c.) Particle 1 : a x (t ) =

dv x
= 2t
dt
dv y
a y (t ) =
=0
dt

Particle 2: a x (t ) =

a x (1.5) = 2(1.5) = 3
a(t ) = 3i.

179

d.) Yes, Particle 2 changes direction at t = 1.7320508 because the x-velocity is zero at
that time.
Total distance =

Total distance =

b
a

1.7320508
0

dy
dx
+ dt
dt
dt

6
t 3 + dt = 4.1592971 4.159.
5
2

Rubric
Part A: 1 point for correctly setting up the differential equation for each particle
1 point for correctly solving the differential equation for each particle
Part B: 1 point for providing the skeleton equation for y, or specifying that one must
solve for an unknown constant
1 point for correctly solving for the unknown constant (either in decimal form or fraction
form)
1 point for providing the correct equation for y
Part C: 1 point for setting up the equation to solve for the acceleration vector for each
particle
1 point for providing the correct acceleration vector for each particle
Part D: 1 point for stating that Particle 2 changes direction at t = 1.7320508
1 point for correctly calculating the total distance traveled to 3 decimal places
3.) a.)

Total= 9 points
x1 = x 0 + x = (1) + (0.1) = 1.1
y 0 = (1) + 2(1)(0.5) 2 = 1.5
y1 = y 0 + xy 0 = (0.5) + (0.1)(1.5) = 0.65
x 2 = x1 + x = (1.1) + (0.1) = 1.2
y1 = (1.1) + 2(1.1)(0.65) 2 = 2.0295
y 2 = y1 + xy1 = (0.65) + (0.1)(2.0295) = 0.85295
x3 = x 2 + x = (1.2) + (0.1) = 1.3
y 2 = (1.2) + 2(1.2)(0.85295) 2 = 2.946056886
y 3 = y 2 + xy 2 = (0.85295) + (0.1)(2.946056886) = 1.147555689

180

x 4 = x3 + x = (1.3) + (0.1) = 1.4

y 3 = (1.3) + 2(1.3)(1.147555689) 2 = 4.723898552
y 4 = y 3 + xy 3 = (1.47555689) + (0.1)(4.723898552) = 1.1947946745

y (1.4) 1.195

b.) While the tables below are not necessary to receive credit, one must show some sort
of calculations that led to the slope field.
dy/dx

(0,0)
0

(0,1)
0

(0,2)
0

(0,3)
0

dy/dx

(1,0)
1

(1,1)
3

(1,2)
9

(1,3)
19

dy/dx

(2,0)
2

(2,1)
6

(2,2)
18

(2,3)
38

dy/dx

(3,0)
3

(3,1)
9

(3,2)
27

(3,3)
57

The slope field should look approximately like this. One is not graded on hand-eye
coordination! As long as the slope field shows a sufficient basis in calculations, full
credit should be received.
c.)

dy
= x + 2xy 2
dx
Polar form: x = r cos , y = r sin

dy
= (r cos ) + 2(r cos )(r sin ) 2
dx

dy
= r cos + 2r 3 sin 2 cos .
dx
Rubric
Part A: 1 point for using the correct formulae for Eulers method
2 points for substituting the correct numerical values for the variables in these equations
1 point for providing the correct approximation to 3 decimal places
181

Part B: 1 point for calculations of values of the slope at each point

2 points for correctly drawing the slope field
Part C: 1 point for showing (in some way) the conversion from Cartesian to polar
coordinates
1 point for correctly representing the derivative in polar form
Total: 9 points
Section II B
4.) a.) The object is gaining altitude on the interval [0, 3]. This is the case because the
derivative of the altitude function (i.e. the graph provided) lies above the horizontal axis
on this interval. When a functions derivative is positive on an interval, that function is
increasing on that same interval.
A
(t ) over a certain time interval gives the total
t
A
accumulated altitude. Thus, summing the areas under the (t ) curve from t = 0 to t = 5
t
yields the total accumulated height:
b.) The definite integral of the function

Note that the 4 came from subtracting the area of the trapezoid from the area of the
neglected triangle that makes up that trapezoid. Summing these areas,
Atotal

0,5

= 9 4 = 5 . Thus, the altitude at t = 5 is 5. Note that one is actually dealing

with an integral here (because the definite integral of a rate over a certain time period
yields a total amount) , so the signs of the areas (i.e. negative and positive) must be
preserved.
c.) This part required a lot of graphical and analytical intuition. In the end, the sketch
should look something like this:

182

Note that there are no tick marks because the problem specified that the sketch need not
be to scale.
Justification: The graph of A/t initially increases linearly, meaning A/t2 (the
derivative) has some constant positive value. A corner then appears, implying a point of
non-differentiability. After this point, A/t decreases linearly, meaning the A/t2 is some
constant negative value. The graph of A/t experiences another corner, meaning it has no
derivative at that point. The graph of A/t then becomes constant for a period of time,
meaning A/t2 is zero for that same time period. After another corner appears, the graph of
the A/t increases linearly, meaning A/t2 is a constant positive value. After another corner
appears, the slope tangent to the semicircle of A/t decreases. Since the semicircle is a
square-root function and the derivative of a square-root function is the reciprocal of a
square-root function, the graph of A/t2 decreases in a reciprocal-square root fashion.
When the graph of A/t reaches a maximum value at the top of the semicircle, the
reciprocal square-root function of A/t2 crosses the t-axis, and continues to decrease. After
one final corner appears, the A/t has a constant value, meaning A/t2 is zero.
b

f ( x)dx

d.) Average value =

. Thus, in order to determine the average value of A/t
a

ba
over the entire interval, one must sum all of the areas (including signs, since one is
dealing with an integral) under the curve and divide by the magnitude of the interval.

While the area of the triangles, trapezoids, and rectangles are simple to calculate, the
semicircular region is a bit troublesome. The area under the curve is not the area of the
semicircle because, since the circular curve is convex rather than concave (i.e. forming a
bowl), the actual area of the semicircle lies below the curve. Thus, one must find the area
of the imaginary rectangle in which the semicircle is inscribed and subtract from this area
the area of the semicircle. After doing all of the geometry, one can sum the areas
(including signs) and divide this by 10 to determine the average value of A/t:

183

Atotal

0 ,10

= 9528+

4=

10.

10
8
= 1.
Average value of A/t =
10
80
Rubric
Part A: 1 point for stating that the object is gaining altitude on [0, 3]
1 point for giving the correct justification
Part B: 1/2 point for summing the areas correctly
1/2 point for supplying the correct total area (total accumulated altitude)
Part C:

2
point for every correctly sketched portion of A/t2 (There are five portions)
5

2
point for every correct justification corresponding with each portion of the graph
5
(There should be a justification for every portion of the graph)
Part D: 1 point for indicating the integral equation for average value
1 point for correctly calculating the average value
Total: 9 points
5.) This is a very conceptual question. This may be good thing for some and signal utter
doom for others (like me the first time I encountered such a question on a test!). It
requires a painfully close analysis of each statement and how the statements relate to each
other. In essence, here is what they mean:
i.) The derivative of f (x) is equal to x 2 g ( x) (simple enough!)
ii.) The derivative of f (x) evaluated at x = 1 is 1 (This is the limit definition of the
derivative)
iii.) As x approaches infinity, f (x) approaches 2.
iv.) f (x) evaluated at 1 is 0 and the second derivative of f (x) evaluated at 1 is 1.
It is now time to synthesize all of this information in each part of the problem:
a.) Since f ( x) = x 2 g ( x),
integral,

(x

(x

g ( x) dx = lim

g ( x) dx = f ( x) + C . Evaluating the improper

(x
a

a 1

g ( x) dx = lim ( f ( x) ) 1 = lim ( f (a ) f (1) ).

a

Statement iii, states that lim f ( x) = 2 and statement iv states that f (1) = 0. Thus,
x

184

lim ( f (a) f (1) ) = 2 0 = 2. Therefore, the improper integral converges to a value of 2.

a

b.) According to statement i, f ( x) = x 2 g ( x) . Furthermore, according to statement ii,

f (1) = 1. Thus, f (1) = (1) 2 g (1) = 1 1 g (1) = 1 g (1) = 0.
c.) lim
x 1

f ( x) f (1) 0
f ( x)
f ( x)
=
= Use LHpitals Rule: lim
= lim
x

1
x

1
g ( x) g (1) 0
g ( x)
g ( x)

f ( x) = x 2 g ( x) f ( x) = 2 x g ( x) g ( x) = 2 x f ( x)
g (1) = 2(1) f (1) = 1
lim
x 1

f ( x) f (1) 1
=
= =1
g ( x) g (1) 1
Rubric
Part A: 1 point for showing that

(x

g ( x) dx = f ( x) + C

1 point for showing that the improper integral converges

1 point for stating the value to which it converges
Part B: 1 point for correctly using the information from statement i
1 point for correctly using the information from statement ii
1 point for showing that g (1) = 0
Part C: 1 point for differentiating the equation is statement i
1 point for finding the correct value of g (1)
f ( x)
= 1.
1 point for showing that lim
x 1 g ( x )

d
6.) a.)
dx

(1) n 2 x n n
2 8 x 18 x 2
= +

+L
(3n + 1)!
4! 7! 10!

d
dx

(1) n 2(0) n n
2
=
(3n + 1)!
4!

n =0

n =0

d2
dx 2

d
dx 2

n =0

n =0

(1) n 2 x n n 8 36 x
=
+L
(3n + 1)!
7! 10!
(1) n 2(0) n n 8
=
(3n + 1)!
7!

185

Since neither the first derivative nor the second derivative is zero when x = 0, there
are no relative extrema or points of inflection at this point.
b.) Max error

1
100

1
f ( n +1) ( )( x c) n +1
. With a range of [-2, 1], f ( n +1) ( ) has the
Rn ( x ) =
100
(n + 1)!
greatest magnitude when it is 1. Thus,
1
(1)(1 0) n +1
1
.

=
100
(n + 1)!
(n + 1)!

1
1
1
1
=
and
=
, which is too large.
(5 + 1)! 120
(4 + 1)! 24
Therefore, one must use a mimimum of a 5th-order Maclaurin polynomial.
c.) Substitute every x with x 3 + 2 :

Using trial and error, n = 5 since

n =0

(1) n 2( x 3 + 2) n n
4th
(3n + 1)!

order f ( x 3 + 2)

2( x 3 + 2) 4( x 3 + 2) 2 6( x 3 + 2) 3 8( x 3 + 2) 4
+

+
4!
7!
10!
13!
Rubric

Part A: 1/2 point for differentiating the power series

1/2 point for finding the correct value of the derivative when x = 0
1/2 for taking the second derivative of the power series
1/2 for finding the correct value of the second derivative when x = 0
1 point for stating that there are no relative extrema or points of inflection
1 point for justification
Part B: 1 point for indicating the use of the Lagrange form of the remainder
1 point for indicating that f ( n +1) ( ) must be 1
1 point for the process of trial and error
1 point for finding the correct order
Part C: 1 point for indicating the substitution to be made
1 point for giving the correct 4th-order Maclaurin polynomial

186

Grading the Practice Exam

A Brief Account on the Theory of AP Grading: Grading AP exams is a science in its
own right. Each year, the final scoring of AP exams is accomplished only after
extremely rigorous statistical analysis. In grading the multiple-choice questions, credit is
awarded based upon each correct answer, and a small fraction is taken away to control for
guessing (the infamous guessing penalty). One is neither rewarded nor penalized for not
answering a question. However, as many have probably heard by now, if even just one
incorrect choice can be ruled out, it is statistically advantageous to guess! To determine
the exact formula used in scoring, AP officials compare the questions on the current
exam with similar similar questions on past exams. By comparing the performance of
students across the country between these two tests, the officials can arrive at an
appropriate formula. During this time, the officials also meet to discuss how the freeresponse sections will be graded. These officials must also agree upon how a composite
score (i.e. the total number of points that a student receives) will be translated into an AP
grade. The AP grades run from scores of 1-5, a 5 signifying extremely well qualified
to receive college credit and a 1 signifying no recommendation for college credit. The
particular formula used to translate a composite score to an AP score varies from year to
year; a certain number of points may equate to a 4 one year and a 3 the next. Not only do
these formulae vary from year to year within a given test, but they also vary significantly
between tests of different subject matter each year. I have friends who could have sworn
that they could expect a 2 on one exam and wound up getting a 4. Similarly, I have also
known people psyched for a 5, only to receive a 4. Indeed, it is quite difficult to
understand the exact mechanics of AP grading. Nevertheless, having knowledge of the
grading process will allow one to strategize accordingly while taking the exam.
Since the formula for grading the AP Calculus BC exam varies from year to
year, one should not expect the one offered in this book to be concrete, for it is very
likely that the formula will change in 2007, 2008, and so on. However, these changes are
not incredibly drastic, so, within certain limits, the AP score received on this practice
exam should be a relatively valid gauge of ones preparedness. The reader should take
note of the questions he or she answered incorrectly or did not understand and review the
concepts associated with these questions. It is not necessarily a good idea to go back and
review material already mastered, as this is often done to stay in ones comfort level. No
matter how painful it is, it is important to tackle that particularly troubling section once
and for all, whatever it may be.
This Books Formula: Step 1: Add up the total number of correctly answered multiplechoice questions.
Step 2: Add up the total number of incorrectly answered multiplechoice questions.
Step 3: To calculate the overall multiple-choice score, multiply the
total number of incorrectly answered multiple-choice questions by 1/4 and subtract this
result from the total number of correctly answered multiple-choice questions.
Step 4: Determine the number points awarded for each of the six
free-response questions.

187

Step 5: Add up the total number of points achieved for the whole
free-response section.
Step 6: Multiply the total multiple-choice score by 1.2. This is the
multiple-choice composite score.
Step 7: Multiply the total free-response score by 1.0. This is the
free-response composite score.
Step 8: Add the results of steps 6 and 7.
Step 9: Find the AP score on the table below.
Composite Score Range
63-108
53-62
42-52
25-41
0-24

AP Grade
5
4
3
2
1

Note that this is my formula based upon how difficult I perceived the pratice exam in
comparison to past exams. This formula will almost certainly differ from the one
actually used.

188

Appendix A: Essential Pre-Calculus Information

Laws of Logarithms
log x a + log x b = log x (ab)
a
log x a log x b = log x
b
log x a y = y log x (a )
log e (a ) ln(a)

( )

log10 (a) = log(a) (When a subscript is not written, the common logarithm is implied.)
log x ( x) = 1 (e.g. ln(e) = 1) (This is the case because the logarithmic function is the inverse
of the exponential function, and vice versa.)
x log x ( a ) = a (For the same reason as above)
Laws of Exponents
(a x )(a y ) = a x + y
(a x )
= a x y
y
(a )
1
a x = x (a 0)
a
Aspects of Square-Root Notation:
n

m
n

a , where a is the radicand, n is the index, and m is the power of the radicand.

Trigonometric Indentities:

sin 2 + cos 2 = 1
tan 2 + 1 = sec 2 (Pythagorean identities)
cot 2 + 1 = csc 2
cos 2 = cos 2 sin 2
sin 2 = 2 sin cos
1 + cos 2
(Double-angle formulae)
cos 2 =
2
1 cos 2
sin 2 =
2
Values of Trigonometric Functions to be Memorized:

189

sin(0) = 0
cos(0) = 1
tan(0) = 0

1
sin =
6 2
3

cos =
2
6
1
tan =
3
6

2

sin =
2
4

2

cos =
2
4

tan = 1
4

sin =
3

cos =
3

tan =
3

3
2
1
2
3

sin = 1
2

cos = 0
2

tan =
2

190

Appendix B: A Brief Table of Integrals

1
du
u
= tan 1 + C
2
a
u +a
a

du = u + C

adu = au + C

du
u 1
2

= sec 1 u + C

(du dv) = du dv = u v + C
u n +1
+ C , (n 1)
n +1

u n du =

du
= ln u + C
u

a
a du =
, a > 0 and a 1

ln a
e u du = e u + C
u

cos

udu

= sin ud

sin udu = cos u + C

tan udu = ln sec u + C
sec udu = ln sec u + tan u + C
csc udu = ln csc u + cot u + C
cot udu = ln sin u + C
191

References
Braden, Bart. Calculating Sums of Infinite Series. The American Mathematical
Monthly, 99(7), Aug. Sep. 1992: 649-655.
Bronson, Richard. Schaums Outline of Theory and Problems of Differential Equations:
Second Edition. New York: McGraw-Hill, 2003
Campbell, Neil A and Jane B. Reece. Biology: Seventh Edition. San Francisco:
Pearson: Benjamin Cummings, 2005.
Chihara, Charles S. On the Possibility of Completing an Infinite Process. The
Philosophical Review, 74(1), Jan 1965: 74-87.
Gray, Alfred. Modern Differential Geometry of Curves and Surfaces with
Mathematica: Second Edition. Boca Raton, FL: CRC Press, 1998.
Kahn, David S. Cracking the AP Calculus AB & BC Exams: 2002-2003 Edition.
New York: Princeton Review Publishing, 2002.
King, Kerry J. and Dale W. Johnson. Cliffs AP Calculus AB and BC: 3rd Edition.
New York: Hungry Minds, 2001.
Metz, Clyde R. Schaums Outline of Theory and Problems of Physical Chemistry:
Second Edition. New York: McGraw-Hill, 1989.
Pelcovits, Robert A. How to Prepare for the AP Physics C Examination. Hauppauge,
NY: Barrons Educational Series, 2002.
Peleg, Yoav; Reuven Pnini; and Elyahu Zaarur. Schuams Outline of Theory and
Problems of Quantum Mechanics. New York: McGraw-Hill, 1998.
Rubinow, S.I. Introduction to Mathematical Biology. New York: Dover, 2002.
Silberberg, Martin S. Chemistry: The Molecular Nature of Matter and Change: Fourth
Edition. New York: McGraw-Hill, 2006.
Tinker, Michael and Robert Lambourne. Further Mathematics for the Physical Sciences.
Chichester, UK: John Wiley & Sons, 2000.
Weir, Maurice D.; Joel Hass; and Frank R. Giordano. Thomas Calculus: Eleventh
Edition. Boston: Pearson: Addison-Wesley, 2005.
Wrede, Robert C. and Murray Spiegel. Schaums Outline of Theory and Problems of
Advanced Calculus: Second Edition. New York: McGraw-Hill, 2002.

192

Index
Absolute convergence, 42
Agnesi, Maria Gaetana, 82
Alternating series, 46, 52-53
Analytic geometry, 75
An Essay on the Principle of Population, 31
Archimedean spiral (see Spirals)
Arc length (in Cartesian coordinates), 75-78
of parametric curves, 84
of polar curves, 100
Astroid, 81
AP Calculus BC, 1
Bernoulli, Jakob, 91
Bounded growth (see Learning curve)
brachistochrone
Carbon-14, 34
Cardioid, 91
Carrying capacity, 28
Cartesian coordinates, 25
Cauchy, Augustin Louis, 66
Cauchy form of the remainder, 66
Cauchys Mean Value Theorem, 8
Center of mass, 7
Circle
parametric, 80
polar, 88
Closed interval, 19
Closed path, 129
Combinatorics, 45
Comparison test (see Infinite series)
Computational cost, 71
Conditional convergence, 42
Conservative field, 130
Conservative force (see Conservative field)
Convergent
improper integrals, 19
infinite series, 40, 42
Cross product (see Vector product)
Curl, 128
Curvature, 123
Del operator, 127
Descartes, Ren, 75

Differential equations, 25-39

solving, 25-27, 72-74
applications of, 27-39
Differential geometry, 122
Directed distance, 87
Directional derivative, 127
Divergence (vector field), 128
Divergent
improper integrals, 19
infinite series, 40
Dot product (see Scalar product)
Drag, 37
Drer, Albrecht, 90
Ellipse (parametric), 80
Euler, Leonhard, 25
Eulers method, 25-27
Field, 126
Field function, 126
Flux, 126
Frenet, Jean Frdric, 124
Frenet frame (see TNB frame)
Free-body diagram, 37
Fundamental Theorem of Algebra, 16
Fundamental Theorem of Calculus, 19
Gabriels horn, 21
Gaussian function, 22
Geometric series, 43
Gradient, 127
Gregory, James, 62
Half-life, 35
Harmonic series, 44
Heisenbergs uncertainty principle, 23
homogeneous (differential equations),
74
Hyperbolic tangent, 39
Ideal gas law, 71
Improper integrals, 19-24
Indeterminate form, 8
193

Infinite sequences, 40-41

Infinite series, 41-54
comparsion test for, 45
integral test for, 44
limit comparison test for, 45
nth term divergence test for, 42
ratio test for, 44
root test for, 45
Initial point (parametric), 79
Initial ray, 87
Integral test (see Infinite series)
Integration
by partial fractions, 15-19
by parts, 11-15
Interval of convergence (see Power series)
Intrinsic rate of increase, 31
Isotherm, 126
Isotope, 34

Madhava of Sangamagrama, 54, 6

Magnetic field, 117
Mathematica, 3, 125
Maxwell, James Clerk, 117
Maxwell-Boltzmann distribution, 23
Methodus incrementum directa
et inversa, 65
Moles, 71
Newton, Isaac, 6, 36
Newtons Law of Cooling, 36
Newtons Second Law of
Motion, 37
Non-conservative force (see Conservative
field)
Normal distribution, 22
Normalized, 22
Open path, 129

Kinetic-molecular theory, 22
Parallelogram method, 106
L'Analyse des Infiniment Petits pour
Parameter, 78
l'Intelligence des Lignes Courbes, 8
Parametric equations and curves, 79-87
Lagrange, Joseph Louis, 66
differentiation of, 82-84
Lagrange form of the remainder, 66
integration of, 84-87
Law of exponential change, 28
Partial derivative, 71, 127
Law of sines and cosines, 106
Partial fraction decomposition, 15
Learning curve, 29
Partial sum, 41
Leibniz, Gottfried, 6
Path independence, 130
Lemniscate, 91-92
Path integral (see Line integral)
Lennard-Jones, John, 69
Pi (case study of), 54
Lennard-Jones potential, 69
Pole, 87
LHpital, Guillaume de, 8
Polar equations and curves, 87-102
LHpitals Rule, 8-11
differentiation of, 94-98
Libby, Willard, 34
integration of, 96-102
Limaon, 90-91
Population ecology, 30
Power series, 56-74
Limit comparison test (see Infinite series)
Linear algebra, 17, 125
functions defined by, 58
Line integral, 128-130
interval of convergence of, 57
Linear restoring force, 68
method of (differential
LIPET, 12
equations), 72
Logistic growth, 28
radius of convergence of, 56
Principal unit normal vector, 123
Logarithmic spiral (see Spirals)
Maclaurin, Colin, 62
Probability density function, 22
Maclaurin polynomial (see Maclaurin series) Products (chemical), 33
Maclaurin series, 62
P-series, 44

194

Quantum mechanics, 23
Radial acceleration, 119
Radioactive decay, 34
Radiometric dating, 34
Radius of convergence (see Power series)
Ramanujan, Srinivasa, 54
Random variable, 22
Rate law, 34
Ratio test, 57 (see also Power series)
Reactants (chemical), 33
Reaction kinetics, 33
Reaction order, 34
Rectangular coordinates (see Cartesian coordinates)
Reduced row echelon form (RREF), 17
Remainder, 65
Repeated factor, 16
Resultant, 104
Right-handed coordinate frame, 111-112
Right hand rule, 114
Rose curve, 92-93
Scalar field, 126-128
Scalar field function (see Field function)
Scalar product, 112-114
Separation of variables, method of, 25
Series (see Infinite series)
Sequences (see Infinite sequences)
Simple harmonic motion, 68-69, 74
Speed, 119
Spirals, 89-90
Step size, 25
Stoichiometric coefficient, 33
Surface area of revolution
of parametric curves, 86
of polar curves, 101
Tangential acceleration, 119
Tautochrone (see Brachistochrone)
Taylor, Brook, 11, 63
Taylor polynomial, 64
Taylor series, 63-74
applications of, 68-74
Taylors theorem, 65
Telescoping series, 44
Terminal point (parametric), 79

Terminal velocity, 39
Term-by-term
differentiation, 59
Term-by-term
integration, 59
TI-89 calculator, 3
TNB frame, 125
Torricelli, Evangelista, 21
Torque, 115-117
Torsion, 125
Transcendental function,
10, 31
Treatise on Fluxions, 62
Truncation, 51-52
Truncation error
(see Truncation)

Underweysung der Messung,

90
Unit vector, 110-111
Unit tangent vector, 123
Unit binormal vector, 123
Vectors (37, 103-131)
addition of, 104-110
calculus of, 118
definition of, 103
multiplication of,
112-117
Vector field, 126-130
Vector function, 118
Vector resolution, 107-110
Vector-valued function
(see Vector function)
Virial equation, 71-72
Wave function, 23
Witch of Agnesi, 82
Work, 113-114, 128
Zenos paradoxes, 5-6, 43

195

196

197

198

199