Professional Documents
Culture Documents
MATHEMATICS 1
WEEK 8: INTEGRATION
2017/2018
CONTENTS
WEEK 8: INTEGRATION............................................................................................................................... 4
8.1 terminology and basic integration rules.................................................................................................. 4
8.2 techniques of integration ....................................................................................................................... 4
8.2.1 integration BY PARTS ......................................................................................................................... 4
8.2.2 Trigonometric SUBSTITUTION ........................................................................................................... 9
8.2.3 PARTIAL FRACTIONS ....................................................................................................................... 11
8.2.3.1 case 1: Q(X) is a product of distinct linear factors ............................................................................ 13
8.2.3.2 case 2: Q(X) is a product of distinct linear factors, some of which are repeated .............................. 14
8.2.3.3 case 3: Q(X) contains irreducible quadratic factors, none of which is repeated................................ 16
8.2.3.4 case 4: Q(X) contains A repeated irreducible quadratic factor..........................................................17
8.3 improper integrals ................................................................................................................................ 18
8.3.1 Type 1: Infinite intervals ..................................................................................................................... 18
8.3.2 Type 2: discontinuous integrands ...................................................................................................... 19
8.3.3 A comparison test for improper integrals........................................................................................... 21
WEEK 9: ENGINEERING APPLICATIONS OF INTEGRALS ......................................................................... 23
9.1 ARC length ........................................................................................................................................... 23
9.1.1 The ARC length function .................................................................................................................... 27
9.2 area of a surface of revolution .............................................................................................................. 28
9.3 hydrostatic pressure and force ............................................................................................................. 34
9.4 moments and centres of mass .............................................................................................................. 36
WEEK 10 : MULTIPLE INTEGRALS ............................................................................................................. 43
10.1 double integrals over rectangles ......................................................................................................... 43
10.1.1 Review of definite integrals .............................................................................................................. 43
10.1.2 volumes and double integrals .......................................................................................................... 43
10.1.3 iterated integrals.............................................................................................................................. 47
10.1.5 average value ................................................................................................................................... 50
10.2 double integrals over general regions ................................................................................................. 51
10.2.1 properties of double integrals .......................................................................................................... 56
10.3 double integrals in polar coordinates .................................................................................................. 57
10.4 triple integrals .................................................................................................................................... 62
1
10.4.1 triple integrals over a general bounded region E .............................................................................. 65
10.4.2 triple integrals in cylindrical coordinates.......................................................................................... 68
10.4.3 cylindrical coordinates ..................................................................................................................... 68
10.4.4 evaluating triple integrals with cylindrical coordinates .................................................................... 70
WEEK 11: DIFFERENTIAL EQUATIONS ..................................................................................................... 73
11.1 INTRODUCTION TO DIFFERENTIAL EQUATION ............................................................................... 73
11.2 THE CLASSIFICATION/TYPE OF DIFFERENTIAL EQUATIONS ........................................................... 75
11.3 SOLUTION TO DIFFERENTIAL EQUATION ........................................................................................ 85
11.4 VERIFICATION OF DIFFERENTIAL EQUATION’S SOLUTION ............................................................. 87
11.5 STRATEGY TO SOLVE 1ST ORDER DIFFERENTIAL EQUATION .......................................................... 91
11.5.1 Exact Differential Equation........................................................................................................... 91
11.5.2 Linear Differential Equation ......................................................................................................... 93
11.5.3 Separable Differential Equation ................................................................................................... 98
11.5.4 Bernoulli’s Differential Equation ................................................................................................ 100
11.5.5 Differential Equation of Homogeneous 𝑑𝑦/𝑑𝑥 = 𝑓(𝑥, 𝑦)/𝑔(𝑥, 𝑦) form ..................................... 105
11.5.6 Differential Equation of Nonhomogeneous 𝑑𝑦/𝑑𝑥 = 𝑓(𝑥, 𝑦)/𝑔(𝑥, 𝑦) Form ..............................107
WEEK 12: SOLUTION TO HOMOGENEOUS & NON-HOMOGENEOUS .................................................. 108
12.1 INTRODUCTION TO 2ND ORDER DIFFERENTIAL EQUATION .......................................................... 108
12.2 SOLUTION TO 2ND ORDER LINEAR DIFFERENTIAL EQUATION ...................................................... 109
12.3 GENERAL THEORY OF LINEARITY PRINCIPLE & LINEAR DEPENDENCY ....................................... 110
12.4 STRATEGY TO SOLVE 2ND ORDER DIFFERENTIAL EQUATION ........................................................115
12.4.1 Homogenous Linear Differential Equation with Constant Coefficients 𝑎, 𝑏, 𝑐 ..............................115
12.4.2 Homogeneous Linear Differential Equation with Non-Constant Coefficients 𝑥2, 𝑎𝑥 (Known as
Euler-Cauchy Differential Equation) ..................................................................................................... 121
12.4.3 Nonhomogeneous Linear Differential Equation with Constant Coefficients 𝑎, 𝑏, 𝑐 In The Form of
𝑟𝑥 = 𝑒𝛼𝑥𝑃𝑛(𝑥).................................................................................................................................... 124
12.4.4 Solving All Types of Nonhomogeneous Linear Differential Equation with Method of Variation of
Parameters .......................................................................................................................................... 142
WEEK 13: POWER SERIES SOLUTIONS FOR DIFFERENTIAL EQUATIONS ..............................................147
13.1 POWER SERIES METHOD .................................................................................................................147
13.1.1 Basic Concepts of Power Series .................................................................................................. 148
13.1.2 Test for Convergence ................................................................................................................. 150
2
13.1.3 Operations of Power Series .........................................................................................................151
13.1.4 Vanishing All Coefficients – A CONDITION THAT IS A BASIC TOOL OF THE POWER SERIES
METHOD ............................................................................................................................................. 152
13.1.5 Idea of The Power Series Method ............................................................................................... 154
13.2 FROBENIUS METHOD ...................................................................................................................... 164
WEEK 14: ENGINEERING APPLICATIONS OF DIFFERENTIAL EQUATION .............................................. 169
14.1 A LIQUID SYSTEM ............................................................................................................................ 169
14.2 A MIXTURE PROBLEM ......................................................................................................................170
14.3 AN LCR CIRCUIT ................................................................................................................................172
14.4 VIBRATING SPRINGS ........................................................................................................................172
14.5 DAMPED VIBRATIONS ......................................................................................................................174
14.6 FORCED VIBRATIONS.......................................................................................................................178
14.7 ELECTRIC CIRCUIT ............................................................................................................................179
Appendix 11.1 MATHEMATICAL MODELING AND ENGINEERING PROBLEM SOLVING ........................ 183
Appendix 11.2 CONSERVATION LAWS AND ENGINEERING ................................................................... 186
3
INTEGRATION
WEEK 8: INTEGRATION
8.1 TERMINOLOGY AND BASIC INTEGRATION RULES
There are two types of integrals: Indefinite and Definite. Indefinite integrals are those with no limits and
definite integrals have limits. When dealing with indefinite integrals, you need to add a constant of
integration. For example, if integrating a function 𝑓(𝑥) with respect to x:
∫ 𝑓(𝑥) 𝑑𝑥 = 𝑔(𝑥) + 𝐶,
Where 𝑔(𝑥) is the integrated function. C is an arbitrary constant called the constant of integration and 𝑑𝑥
indicated the variable with respect to which we are integrating, in this case,𝑥. The function being integrated,
𝑓(𝑥), is called the integrand.
The integral of many functions are well known, and there are useful rules to work out the integral of more
complication functions, which are shown in Figure 8.1 below. The summary of the common procedures for
fitting integrands to the basic integration rules is given in Figure 8.2.
This section will discussed in more detail three methods of integration: Integration by parts, the substitution
method and partial fractions.
One of the important integration techniques is called integration by parts. This technique can be applied to
a wide variety of functions and is particularly useful for integrands involving products of algebraic and
transcendental functions.
4
Figure 8.1: Basic Integration Rules (𝑎 > 0)
5
Figure 8.2: Procedures for Fitting Integrands to Basic Rule
Let 𝑢 = 𝑓(𝑥) and 𝑣 = 𝑔(𝑥). Then the differentials are 𝑑𝑢 = 𝑓 ′ (𝑥) 𝑑𝑥 and 𝑑𝑣 = 𝑔′ (𝑥) 𝑑𝑥. So, by the
Substitution Rule, the formula for integration by parts becomes
∫ 𝑢 𝑑𝑣 = 𝑢𝑣 − ∫ 𝑣 𝑑𝑢
Example
1. Find ∫ 𝑥 sin 𝑥 𝑑𝑥
Solution
First method
Suppose we choose 𝑓(𝑥) = 𝑥 and 𝑔′ (𝑥) = sin 𝑥. Then 𝑓 ′ (𝑥) = 1 and 𝑔(𝑥) = − cos 𝑥. Note that for 𝑔, we
can choose any antiderivative of 𝑔′ . Thus, using the formula in Eq 2,
= −𝑥 cos 𝑥 + ∫ cos 𝑥 𝑑𝑥
= −𝑥 cos 𝑥 + sin 𝑥 + 𝐶
6
Second method
Let:
𝑢=𝑥 𝑑𝑣 = sin 𝑥 𝑑𝑥
𝑑𝑢 = 𝑑𝑥 𝑣 = − cos 𝑥
u dv
∫ 𝑥 sin 𝑥 𝑑𝑥 = ∫ 𝑥 sin 𝑥 𝑑𝑥
u v u du
= −𝑥 cos 𝑥 + ∫ cos 𝑥 𝑑𝑥
= −𝑥 cos 𝑥 + sin 𝑥 + 𝐶
We can evaluate definite integrals by parts. By evaluating both sides of Eq 2 (formula for integration by
parts) between 𝑎 and 𝑏, assuming 𝑓 ′ and 𝑔′ are continuous, and using the Fundamental Theorem of
Calculus, we get
𝑏 𝑏 𝑏
′ (𝑥)𝑑𝑥
∫ 𝑓(𝑥)𝑔 = 𝑓(𝑥)𝑔(𝑥)] − ∫ 𝑔(𝑥) 𝑓 ′ (𝑥) 𝑑𝑥
𝑎 𝑎 𝑎
We could also use trigonometric identities to integrate certain combinations of trigonometric functions.
Example
Solution
We could convert cos 2 𝑥 to 1 − sin2 𝑥, but we would be left with an expression in terms of sin 𝑥 with no
extra cos 𝑥 factor. Therefore, we could separate a single sine factor and rewrite the remaining sin4 factor
in terms of cos 𝑥 factor.
𝑢3 𝑢5 𝑢7
= −( − 2 + )+ 𝐶
3 5 7
1 2 1
= − cos 3 𝑥 + cos 5 𝑥 − cos7 𝑥 + 𝐶
3 5 7
We can use a similar strategy to evaluate integrals of the form ∫ tan𝑚 𝑥 sec 𝑛 𝑥 𝑑𝑥, where 𝑚 > 0, 𝑛 > 0 are
integers.
𝑑
Since (𝑑𝑥) tan 𝑥 = sec 2 𝑥, we could separate a sec 2 𝑥 factor and convert the remaining (even) power of
secant to an expression involving tangent using the following identity:
sec 2 𝑥 = 1 + tan2 𝑥
Another way is to separate a sec 𝑥 tan 𝑥 factor and convert the remaining (even) power of tangent to
𝑑
secant since ( ) sec 𝑥 = sec 𝑥 tan 𝑥 .
𝑑𝑥
As for other cases, the use of identities, integration by parts, and occasionally a little ingenuity may come
handy. The following formulas and trigonometric identities are also useful:
∫ tan 𝑥 𝑑𝑥 = ln|sec 𝑥| + 𝐶
1
sin 𝐴 cos 𝐵 = [sin(𝐴 − 𝐵) + sin(𝐴 + 𝐵)]
2
1
sin 𝐴 sin 𝐵 = [cos(𝐴 − 𝐵) − cos(𝐴 + 𝐵)]
2
1
cos 𝐴 cos 𝐵 = [cos(𝐴 − 𝐵) + cos(𝐴 + 𝐵)]
2
8
8.2.2 TRIGONOMETRIC SUBSTITUTION
In finding the area of a circle or an ellipse, an integral of the form ∫ √𝑎2 − 𝑥 2 𝑑𝑥 arises, where 𝑎 > 0. If it
were ∫ 𝑥√𝑎2 − 𝑥 2 𝑑𝑥, the substitution 𝑢 = 𝑎2 − 𝑥 2 would be effective. However,
∫ √𝑎2 − 𝑥 2 𝑑𝑥 … … (3)
Eq (3) would be more challenging. If we change the variable from 𝑥 to 𝜃 by the substitution of 𝑥 = 𝑎 sin 𝜃.
Then, the root sign of Eq (3) can be removed by making use of the identity 1 − sin2 𝜃 = cos2 𝜃. This is
shown as below:
= √𝑎2 cos2 𝜃
= 𝑎| cos 𝜃 |
Notice the difference between the substitution 𝑢 = 𝑎2 − 𝑥 2 (in which the new variable is a function of the
old one) and the substitution 𝑥 = 𝑎 sin 𝜃 (the old variable is a function of the new one).
In general, we can make a substitution of the form 𝑥 = 𝑔(𝑡) by using the Substitution Rule in reverse. To
make our calculations simpler, we assume that g has an inverse function; that is, g is one-to-one.
This type of substitution is called the inverse substitution. The inverse substitution of 𝑥 = 𝑎 sin 𝜃 can be
made provided that it defines a one-to-one function. This can be accomplished by restricting 𝜃 to lie in the
𝜋 𝜋
interval [− , ]
2 2
Table 8.1 below shows a list of trigonometric substitutions which are effective for the given radical
expressions because of the specified trigonometric identities. The restriction on 𝜃 is imposed in each of the
cases shown in Table 8.1 to ensure that the function that defines the substitution is one-to-one.
9
Table 8.1: Table of Trigonometric Substitution
Example
3. Evaluate
√9 − 𝑥 2
∫ 𝑑𝑥
𝑥2
Solution
Let 𝑥 = 3 sin 𝜃, where −𝜋/2 ≤ 𝜃 ≤ 𝜋/2. Then, 𝑑𝑥 = 3 cos 𝜃 𝑑𝜃.
√9 − 𝑥 2 = √9 − 9sin2 𝜃
= √9 cos 2 𝜃
= 3| cos 𝜃|
= 3 cos 𝜃
(Note that cos 𝜃 ≥ 0 because −𝜋/2 ≤ 𝜃 ≤ 𝜋/2.) Thus using the Inverse Substitution Rule:
√9 − 𝑥 2 3 cos 𝜃
∫ 2
𝑑𝑥 = ∫ 3 cos 𝜃 𝑑𝜃
𝑥 9sin2 𝜃
cos2 𝜃
=∫ 𝑑𝜃
sin2 𝜃
= ∫ cot 2 𝜃 𝑑𝜃
= ∫(csc 2 𝜃 − 1)𝑑𝜃
= − cot 𝜃 − 𝜃 + 𝐶
10
Since this is an indefinite integral, we must return to the original variable x. This can be done either by
using trigonometric identities to express cot 𝜃 in terms of sin 𝜃 = 𝑥/3 or by drawing a diagram, as in Figure
8.3, where 𝜃 is interpreted as an angle of a right triangle.
Based on the Pythagorean Theorem, the length of the adjacent side can be expressed as
√9 − 𝑥 2
Then, we can simply read the value of cot 𝜃 from the figure:
√9 − 𝑥 2
cot 𝜃 =
𝑥
(Although 𝜃 > 0 in the diagram, this expression for cot 𝜃 is valid even when 𝜃 < 0). Since sin 𝜃 = 𝑥/3, then
𝜃 = sin−1 (𝑥/3). Therefore
√9 − 𝑥 2 √9 − 𝑥 2 −1
𝑥
∫ 𝑑𝑥 = − − sin ( )+𝐶
𝑥2 𝑥 3
This section demonstrates a method to integrate any rational function (a ratio of polynomials) by expressing
it as a sum of simpler fractions, called partial fractions, that we already know how to integrate.
2 1
To illustrate the method, observe that by taking the fractions and to a common denominator,
(𝑥−1) (𝑥+2)
the expression becomes
2 1 2(𝑥 + 2) − (𝑥 − 1)
− =
(𝑥 − 1) (𝑥 + 2) (𝑥 − 1)(𝑥 + 2)
𝑥+5
=
𝑥2 +𝑥−2
If we reverse the procedure, we see how to integrate the function on the right side of the following equation.
11
𝑥+5 2 1
∫ 𝑑𝑥 = ∫ ( − ) 𝑑𝑥
𝑥2 +𝑥−2 (𝑥 − 1) (𝑥 + 2)
= 2 ln|𝑥 − 1| − ln|𝑥 + 2| + 𝐶
𝑃(𝑥)
𝑓(𝑥) =
𝑄(𝑥)
Be a rational function where 𝑃(𝑥) and 𝑄(𝑥) are polynomials. The function 𝑓(𝑥) can be expressed as a sum
of simpler fractions provided that the degree of P is less than the degree of Q. Such a rational function is
called proper.
If
Where 𝑎𝑛 ≠ 0, then the degree of P is n and we write deg(P) = n. If f is improper, that is, deg(P) deg(Q),
then we must take the preliminary step of dividing Q into P (by long division) until a remainder R (x) is
obtained such that deg(R) < deg(Q).
𝑃(𝑥) 𝑅(𝑥)
𝑓(𝑥) = = 𝑆(𝑥) + … . . (4)
𝑄(𝑥) 𝑄(𝑥)
As the next example illustrates, sometimes this preliminary step is all that is required.
Example
4. Find
𝑥3 + 𝑥
∫ 𝑑𝑥
𝑥−1
Solution
Since the degree of the numerator is greater than the degree of the denominator, we first perform the long
division.
𝑥3 + 𝑥 2
∫ 𝑑𝑥 = ∫ (𝑥 2 + 𝑥 + 2 + ) 𝑑𝑥
𝑥−1 𝑥−1
𝑥3 𝑥2
= + + 2𝑥 + 2 ln|𝑥 − 1| + 𝐶
3 2
12
From Eq (4), if the denominator is more complicated, then the next step is to factor the denominator Q
(x) as far as possible. It can be shown that any polynomial Q can be factored as a product of linear factors
(of the form ax + b) and irreducible quadratic factors (of the form 𝑎𝑥 2 + 𝑏𝑥 + 𝑐, where 𝑏 2 − 4𝑎𝑐 < 0.
Then, the next step is to express the proper rational function 𝑅(𝑥)/𝑄(𝑥) in Eq (4) as a sum of partial
fractions of the following form
𝐴 𝐴𝑥 + 𝐵
𝑜𝑟
(𝑎𝑥 + 𝑏)𝑖 (𝑎𝑥 2 + 𝑏𝑥 + 𝑐)𝑗
A theorem in algebra guarantees that it is always possible to do this. We explain the details for the four
cases that occur:
Where no factor is repeated and no factor is a constant multiple of another. Hence, in this case, the partial
fraction theorem states that there exist constant 𝐴1 , 𝐴2 , … , 𝐴𝑘 such that
𝑅(𝑥) 𝐴1 𝐴2 𝐴𝑘
= + +⋯+ … . (5)
𝑄(𝑥) 𝑎1 𝑥 + 𝑏1 𝑎2 𝑥 + 𝑏2 𝑎𝑘 𝑥 + 𝑏𝑘
Example
5. Evaluate
𝑥 2 + 2𝑥 − 1
∫ 𝑑𝑥
2𝑥 3 + 3𝑥 2 − 2𝑥
Solution
Since the degree of the numerator is less than the degree of the denominator, we don’t need to divide. We
factor the denominator as
2𝑥 3 + 3𝑥 2 − 2𝑥 = 𝑥(2𝑥 2 + 3𝑥 − 2)
= 𝑥(2𝑥 − 1)(𝑥 + 2)
13
Then,
𝑥 2 + 2𝑥 − 1 𝐴 𝐵 𝐶
= + +
𝑥(2𝑥 − 1)(𝑥 + 2) 𝑥 2𝑥 − 1 𝑥 + 2
In order to determine the constant 𝐴, 𝐵 and C, multiply both sides of the equation by the product of the
denominators to give
System of equations:
2𝐴 + 𝐵 + 2𝐶 = 1
3𝐴 + 2𝐵 − 𝐶 = 2
−2𝐴 = −1
1 1 1
Thus, 𝐴 = , 𝐵 = , 𝐶 = −
2 5 10
𝑥 2 + 2𝑥 − 1 11 1 1 1 1
∫ 3 2
𝑑𝑥 = ∫ ( + − ) 𝑑𝑥
2𝑥 + 3𝑥 − 2𝑥 2 𝑥 5 2𝑥 − 1 10 𝑥 + 2
1 1 1
= ln|𝑥| + ln|2𝑥 − 1| − ln|𝑥 + 2| + 𝐾
2 10 10
Note that in integrating the middle term, the following substitutions have been made:
1
𝑢 = 2𝑥 − 1, 𝑑𝑢 = 2𝑑𝑥, then 𝑑𝑥 = 𝑑𝑢
2
8.2.3.2 CASE 2: Q(X) IS A PRODUCT OF DISTINCT LINEAR FACTORS, SOME OF WHICH ARE
REPEATED
Suppose the first linear factor (𝑎1 𝑥 + 𝑏1 ) is repeated 𝑟 times, that is, (𝑎1 𝑥 + 𝑏1 )𝑟 occurs in the factorization
of 𝑄(𝑥). Then, instead of the single term in Eq (5), we could use
𝐴1 𝐴2 𝐴𝑟
+ 2
+⋯+ … . (6)
𝑎1 𝑥 + 𝑏1 (𝑎1 𝑥 + 𝑏1 ) (𝑎1 𝑥 + 𝑏1 )𝑟
14
Example
6. Find
𝑥 4 − 2𝑥 2 + 4𝑥 + 1
∫ 𝑑𝑥
𝑥3 − 𝑥2 − 𝑥 + 1
Solution
𝑥 4 − 2𝑥 2 + 4𝑥 + 1 4𝑥
3 2
=𝑥+1+ 3 2
𝑥 −𝑥 −𝑥+1 𝑥 −𝑥 −𝑥+1
𝑥 3 − 𝑥 2 − 𝑥 + 1 = (𝑥 − 1)(𝑥 2 − 1)
= (𝑥 − 1)(𝑥 − 1)(𝑥 + 1)
= (𝑥 − 1)2 (𝑥 + 1)
Hence,
4𝑥 𝐴 𝐵 𝐶
2
= + 2
+
(𝑥 − 1) (𝑥 + 1) 𝑥 − 1 (𝑥 − 1) 𝑥+1
𝑥 4 − 2𝑥 2 + 4𝑥 + 1 1 2 1
∫ 3 2
𝑑𝑥 = ∫ [𝑥 + 1 + + 2
− ] 𝑑𝑥
𝑥 −𝑥 −𝑥+1 𝑥 − 1 (𝑥 − 1) 𝑥+1
𝑥2 2
= + 𝑥 + ln|𝑥 − 1| − − ln|𝑥 + 1| + 𝐾
2 𝑥−1
𝑥2 2 𝑥−1
= +𝑥− + ln | | +𝐾
2 𝑥−1 𝑥+1
15
8.2.3.3 CASE 3: Q(X) CONTAINS IRREDUCIBLE QUADRATIC FACTORS, NONE OF WHICH IS
REPEATED
If 𝑄(𝑥) has the factor 𝑎𝑥 2 + 𝑏𝑥 + 𝑐, where 𝑏 2 − 4𝑎𝑐 < 0, then, in addition to the partial fractions in Eq (5)
and (6), the expression for 𝑅(𝑥)/𝑄(𝑥) will have a term of the form
𝐴𝑥 + 𝐵
… . . (7)
𝑎𝑥 2 + 𝑏𝑥 + 𝑐
Where A and B are constants to be determined. The term in Eq (7) can be integrated by completing squares
(if necessary) and using the following formula
𝑑𝑥 1 𝑥
∫ = tan−1 ( ) + 𝐶
𝑥2 +𝑎 2 𝑎 𝑎
Example
7. Evaluate
4𝑥 2 − 3𝑥 + 2
∫ 𝑑𝑥
4𝑥 2 − 4𝑥 + 3
Solution
Since the degree of the numerator is not less than the degree of the denominator, we divide the
expression, which yield
4𝑥 2 − 3𝑥 + 2 𝑥−1
2
=1+ 2
4𝑥 − 4𝑥 + 3 4𝑥 − 4𝑥 + 3
Note that the quadratic 4𝑥 2 − 4𝑥 + 3 is irreducible because its discriminant 𝑏 2 − 4𝑎𝑐 = −32 < 0. Hence,
we complete the square
4𝑥 2 − 4𝑥 + 3 = (2𝑥 − 1)2 + 2
Let
1
𝑢 = 2𝑥 − 1, 𝑑𝑢 = 2 𝑑𝑥, 𝑥 = (𝑢 + 1)
2
Hence,
4𝑥 2 − 3𝑥 + 2 𝑥−1
∫ 2
𝑑𝑥 = ∫ (1 + 2 ) 𝑑𝑥
4𝑥 − 4𝑥 + 3 4𝑥 − 4𝑥 + 3
1
1 2 (𝑢 + 1) − 1
=𝑥+ ∫ 𝑑𝑢
2 𝑢2 + 2
1 𝑢−1
=𝑥+ ∫ 2 𝑑𝑢
4 𝑢 +2
16
1 𝑢 1 1
=𝑥+ ∫ 2 𝑑𝑢 − ∫ 2 𝑑𝑢
4 𝑢 +2 4 𝑢 +2
1 1 1 𝑢
= 𝑥 + ln(𝑢2 + 2) − . tan−1 ( ) + 𝐶
8 4 √2 √2
1 1 2𝑥 − 1
= 𝑥 + ln(4𝑥 2 − 4𝑥 + 3) − tan−1 ( )+𝐶
8 4√2 √2
If 𝑄(𝑥) has the factor (𝑎𝑥 2 + 𝑏𝑥 + 𝑐)𝑟 , where 𝑏 2 − 4𝑎𝑐 < 0, then
𝐴1 𝑥 + 𝐵1 𝐴2 𝑥 + 𝐵2 𝐴𝑟 𝑥 + 𝐵𝑟
+ + ⋯ … . . (8)
𝑎𝑥 2 + 𝑏𝑥 + 𝑐 (𝑎𝑥 2 + 𝑏𝑥 + 𝑐)2 (𝑎𝑥 2 + 𝑏𝑥 + 𝑐)𝑟
occurs in the partial fraction decomposition of 𝑅(𝑥)/𝑄(𝑥). Each of the terms in Eq (8) can be integrated
using a substitution or by first completing the square if necessary.
Example
8. Evaluate
1 − 𝑥 + 2𝑥 2 − 𝑥 3
∫ 𝑑𝑥
𝑥(𝑥 2 + 1)2
Solution
1 − 𝑥 + 2𝑥 2 − 𝑥 3 𝐴 𝐵𝑥 + 𝐶 𝐷𝑥 + 𝐸
2 2
= + 2 + 2
𝑥(𝑥 + 1) 𝑥 𝑥 + 1 (𝑥 + 1)2
Then
1 − 𝑥 + 2𝑥 2 − 𝑥 3 1 𝑥+1 𝑥
∫ 2 2
𝑑𝑥 = ∫ ( − 2 + 2 ) 𝑑𝑥
𝑥(𝑥 + 1) 𝑥 𝑥 + 1 (𝑥 + 1)2
𝑑𝑥 𝑥 𝑑𝑥 𝑥 𝑑𝑥
=∫ −∫ 2 𝑑𝑥 − ∫ 2 +∫ 2
𝑥 𝑥 +1 𝑥 +1 (𝑥 + 1)2
1 1
= ln|𝑥| − ln(𝑥 2 + 1) − tan−1 𝑥 − 2
+𝐾
2 2(𝑥 + 1)
17
8.3 IMPROPER INTEGRALS
In this sub section, extend the concept of a definite integral to the case where the interval is infinite and
also to the case where f has an infinite discontinuity in [a, b]. In either case the integral is called an
improper integral
Consider the infinite region 𝒮 that lies under the curve 𝑦 = 1/𝑥 2 , above the x-axis, and to the right of line
𝑥 = 1. This is shown in Figure 8.4.
1
=1−
𝑡
Notice that 𝐴(𝑡) < 1 no matter how large 𝑡 is chosen. We also observe that
1
lim 𝐴(𝑡) = lim (1 − ) = 1
𝑡→∞ 𝑡→∞ 𝑡
This is shown in Figure 8.5.
Example
∞ 1
9. Determine whether the integral ∫1 (𝑥) 𝑑𝑥 is convergent or divergent
Solution
∞
1
∫ ( ) 𝑑𝑥 = lim ln|𝑥|]1𝑡 = lim (ln 𝑡 − ln 1)
1 𝑥 𝑡→∞ 𝑡→∞
= lim ln 𝑡 = ∞
𝑡→∞
The limit does not exist as a finite number and so the improper integral is divergent.
Suppose that 𝑓 is a positive continuous function defined on a finite interval [a, b) but has a vertical asymptote
at b. Let S be the unbounded region under the graph of f and above the x-axis between a and b. (For Type 1
integrals, the regions extended indefinitely in a horizontal direction. Here the region is infinite in a vertical
direction.)
19
The area of the part of S between a and t (the shaded region in Figure 8.7) is
𝑡
𝐴(𝑡) = ∫ 𝑓(𝑥)𝑑𝑥
𝑎
If it happen that 𝐴(𝑡) approaches a definite number 𝐴 as 𝑡 → 𝑏 − , then we say that the area of the region S
is 𝐴 and we could write
𝑏 𝑏
∫ 𝑓(𝑥)𝑑𝑥 = lim− ∫ 𝑓(𝑥)𝑑𝑥
𝑎 𝑡→𝑏 𝑎
We use this equation to define an improper integral of Type 2 even when 𝑓 is not a positive function, no
matter what type of discontinuity 𝑓 has at b. Figure 8.8 below presents the definition of an improper
Integral of Type 2.
10. Find
5
1
∫ 𝑑𝑥
2 √𝑥 − 2
Solution
1
The given integral is improper because 𝑓(𝑥) = 𝑥−2
has the vertical asymptote at 𝑥 = 2. Since the infinite
√
discontinuity occurs at the left endpoint of [2, 5], we use part (b) of the definition in Figure 8.8.
5 5
1 𝑑𝑥
∫ 𝑑𝑥 = lim+ ∫
2 √𝑥 − 2 𝑡→2 𝑡 √𝑥 − 2
5
= lim+2√𝑥 − 2]𝑡 = lim+(√3 − √𝑡 − 2)
𝑡→2 𝑡→2
= 2√3
Thus the given improper integral is convergent and, since the integrand is positive, we can interpret the
value of the integral as the area of the shaded region in Figure 8.9.
1
Figure 8.9: 𝑦 =
√𝑥−2
Sometimes it is impossible to find the exact value of an improper integral and yet it is important to know
whether it is convergent or divergent.
In such cases the following theorem is useful. Although we state it for Type 1 integrals, a similar theorem is
true for Type 2 integrals.
21
We omit the proof of the Comparison Theorem, but Figure 8.10 makes it seem plausible.
If the area under the top curve y = f (x) is finite, then so is the area under the bottom curve y = g (x).
If the area under y = g (x) is infinite, then so is the area under y = f (x). [Note that the reverse is not
∞ ∞ ∞
necessarily true: If ∫𝑎 𝑔(𝑥)𝑑𝑥 is convergent, ∫𝑎 𝑓(𝑥)𝑑𝑥 may or may not be convergent, and if ∫𝑎 𝑓(𝑥)𝑑𝑥
∞
is divergent, ∫𝑎 𝑔(𝑥)𝑑𝑥 may or may not be divergent.]
22
ENGINEERING APPLICATIONS OF
INTEGRALS
WEEK 9: ENGINEERING APPLICATIONS OF INTEGRALS
9.1 ARC LENGTH
What do we mean by the length of a curve? We might think of fitting a piece of string to the curve in Figure
9.1 and then measuring the string against a ruler. But that might be difficult to do with much accuracy if we
have a complicated curve.
We need a precise definition for the length of an arc of a curve, in the same spirit as the definitions we
developed for the concepts of area and volume.
If the curve is a polygon, we can easily find its length; we just add the lengths of the line segments that form
the polygon. (We can use the distance formula to find the distance between the endpoints of each
segment).
We are going to define the length of a general curve by first approximating it by a polygon and then taking a
limit as the number of segments of the polygon is increased. This process is familiar for the case of a circle,
where the circumference is the limit of lengths of inscribed polygons (see Figure 9.2).
Suppose that a curve 𝐶 is defined by the equation 𝑦 = 𝑓(𝑥), where 𝑓 is continuous and 𝑎 ≤ 𝑥 ≤ 𝑏. We
obtain a polygonal approximation to 𝐶 by dividing the interval [a, b] into n subintervals with endpoints
𝑥0 , 𝑥1 , 𝑥2 , … . , 𝑥𝑛 and equal width ∆𝑥.
If 𝑦𝑖 = 𝑓(𝑥𝑖 ), then the point 𝑃𝑖 (𝑥𝑖 𝑦𝑖 ) lies on 𝐶 and the polygon with vertices 𝑃0 , 𝑃1 , … . , 𝑃𝑛 illustrated in
Figure 9.3, is an approximation to 𝐶
23
Figure 9.3: Polygon with vertices 𝑃0 , 𝑃1 , … . , 𝑃𝑛
The length 𝐿 of 𝐶 is approximately the length of this polygon and the approximation gets better as we let n
increase. (See Figure 9.4, where the arc of the curve between 𝑃𝑖−1 and 𝑃𝑖 has been magnified and
approximations with successively smaller values of ∆𝑥 are shown.)
Figure 9.4: Arc of the curve between 𝑃𝑖−1 and 𝑃𝑖 has been magnified
Therefore we define the length 𝐿 of the curve 𝐶 with equation, 𝑦 = 𝑓(𝑥), 𝑎 ≤ 𝑥 ≤ 𝑏 as the limit of the
lengths of these inscribed polygons ( if the limit exists):
𝑛
Note that the procedure for defining arc length is very similar to the procedure we used for defining area
and volume: We divided the curve into a large number of small parts. We then found the approximate
lengths of the small parts and added them. Finally, we took the limit as 𝑛 → ∞.
The definition of arc length given by Eq. (1) is not very convenient for computational purposes, but we can
derive an integral formula for 𝐿 in the case where 𝑓 has a continuous derivative. [Such a function 𝑓 is called
smooth because a small change in 𝑥 produces a small change in 𝑓 ′ (𝑥).]
24
Let ∆𝑦𝑖 = 𝑦𝑖 − 𝑦𝑖−1 . Then
Applying the Mean Value Theorem of 𝑓 on the interval [𝑥𝑖−1 , 𝑥𝑖 ], we find that there is a number 𝑥𝑖∗
between 𝑥𝑖−1 and 𝑥𝑖 such that
Thus, we have
|𝑃𝑖−1 𝑃𝑖 | = √(∆𝑥)2 + (∆𝑦𝑖 )2 = √(∆𝑥)2 + [𝑓 ′ (𝑥𝑖∗ )∆𝑥]2
by the definition of a definite integral. We know that this integral exists because the function 𝑔(𝑥) =
√1 + [𝑓 ′ (𝑥)]2 is continuous. Thus we have proved the following theorem:
If we use Leibniz notation for derivatives, we can write the arc length formula as follows:
𝑏
𝑑𝑦 2
𝐿 = ∫ √1 + ( ) 𝑑𝑥 … … (3)
𝑎 𝑑𝑥
25
Example
1. Find the length of the arc of the semicubical parabola 𝑦 2 = 𝑥 3 between the points (1, 1) and (4, 8).
See Figure 9.5
Figure 9.5: 𝑦 2 = 𝑥 3
Solution
4
𝑑𝑦 2 4
9
√
𝐿 = ∫ 1 + ( ) 𝑑𝑥 = ∫ √1 + 𝑥 𝑑𝑥
1 𝑑𝑥 1 4
9 9
If we substitute 𝑢 = 1 + 4 𝑥 , then 𝑑𝑢 = 4 𝑑𝑥
13
When x = 1, 𝑢 = 4
; when x = 4, u = 10. Therefore
10
4 10 4 2
𝐿 = ∫ √𝑢 𝑑𝑢 = ∙ 𝑢3/2 |
9 13/4 9 3 13/4
8 13 3/2
= [103/2 − ( ) ]
27 4
1
= (80√10 − 13√13)
27
If a curve has the equation x = g (y), c y d, and g (y) is continuous, then by interchanging the roles of x
and y in Formula 2 or Eq. (3), we obtain the following formula for its length:
26
9.1.1 THE ARC LENGTH FUNCTION
We will find it useful to have a function that measures the arc length of a curve from a particular starting
point to any other point on the curve.
Thus if a smooth curve C has the equation, y = f (x), a x b let s (x) be the distance along C from the initial
point P0(a, f (a)) to the point Q (x, f (x)). Then s is a function, called the arc length function, and, by Formula
2,
𝑥
𝑠(𝑥) = ∫ √1 + [𝑓 ′ (𝑡)]2 𝑑𝑡 … (5)
𝑎
(We have replaced the variable of integration by t so that x does not have two meanings.) We can use the
Fundamental Theorem of Calculus to differentiate Eq. (5) (since the integrand is continuous):
𝑑𝑠 𝑑𝑦 2
= √1 + [𝑓 ′ (𝑥)]2 = √1 + ( ) … (6)
𝑑𝑥 𝑑𝑥
Eq. (6) shows that the rate of change of s with respect to x is always at least 1 and is equal to 1 when f(x),
the slope of the curve, is 0. The differential of arc length is
𝑑𝑦 2
𝑑𝑠 = √1 + ( ) 𝑑𝑥 … (7)
𝑑𝑥
The geometric interpretation of Eq. (8) is shown in Figure 9.6. It can be used as a mnemonic
device for remembering both of the Formulas 3 and 4.
If we write 𝐿 = ∫ 𝑑𝑠, then from Eq. (8) either we can solve to get (7), which gives (3), or we can solve to get
27
𝑑𝑥 2
𝑑𝑠 = √1 + ( ) 𝑑𝑦
𝑑𝑦
Example
Find the arc length function for the curve 𝑦 = 𝑥 2 − 1⁄8 ln 𝑥 taking P0(1, 1) as the starting point.
Solution
1
If 𝑓(𝑥) = 𝑥 2 − 1⁄8 ln 𝑥 , then 𝑓 ′ (𝑥) = 2𝑥 − 8𝑥
1 2 1 1
1 + [𝑓 ′ (𝑥)]2 = 1 + (2𝑥 − ) = 1 + 4𝑥 2 − +
8𝑥 2 64𝑥 2
1 1 1 2
= 4𝑥 2 + + = (2𝑥 + )
2 64𝑥 2 8𝑥
1
√1 + [𝑓 ′ (𝑥)]2 = 2𝑥 + ,
8𝑥
Since 𝑥 > 0. Thus the arc length function is given by
𝑥
𝑠(𝑥) = ∫ √1 + [𝑓 ′ (𝑡)]2 𝑑𝑡
1
𝑥
1 𝑥
= ∫ (2𝑡 + ) 𝑑𝑡 = 𝑡 2 + 1⁄8 ln 𝑡|
1 8𝑡 1
= 𝑥 2 + 1⁄8 ln 𝑥 − 1
For instance, the arc length along the curve from (1, 1) to (3, f (3)) is
ln 3
𝑠(3) = 32 + 1⁄8 ln 3 − 1 = 8 + ≈ 8.1373
8
A surface of revolution is formed when a curve is rotated about a line. Such a surface is the lateral boundary
of a solid of revolution.
We want to define the area of a surface of revolution in such a way that it corresponds to our intuition.
If the surface area is A, we can imagine that painting the surface would require the same amount of paint
as does a flat region with area A.
28
Let’s start with some simple surfaces. The lateral surface area of a circular cylinder with radius r and height
h is taken to be A = 2 rh because we can imagine cutting the cylinder and unrolling it (as in Figure 9.7) to
obtain a rectangle with dimensions 2 r and h.
Likewise, we can take a circular cone with base radius r and slant height l, cut it along the dashed line in
2𝜋𝑟
Figure 9.8, and flatten it to form a sector of a circle with radius l and central angle 𝜃 = 𝑙
.
1 2
In general, the area of a sector of a circle with radius l and angle is 2
𝑙 𝜃 and so in this case the area is
1 1 2𝜋𝑟
𝐴 = 𝑙2 𝜃 = 𝑙2 ( ) = 𝜋𝑟𝑙
2 2 𝑙
Therefore we define the lateral surface area of a cone to be 𝐴 = 𝜋𝑟𝑙. What about more complicated
surfaces of revolution? If we follow the strategy we used with arc length, we can approximate the original
curve by a polygon. When this polygon is rotated about an axis, it creates a simpler surface whose surface
area approximates the actual surface area.
By taking a limit, we can determine the exact surface area. The approximating surface, then, consists of a
number of bands, each formed by rotating a line segment about an axis. To find the surface area, each of
these bands can be considered a portion of a circular cone, as shown in Figure 9.9.
29
Figure 9.8: A portion of circular cone
The area of the band (or frustum of a cone) with slant height l and upper and lower radii r1 and r2 is found
by subtracting the areas of two cones:
𝑙1 𝑙1 + 𝑙
=
𝑟1 𝑟2
Which gives
𝑟2 𝑙1 = 𝑟1 𝑙1 + 𝑟1 𝑙 or (𝑟2 − 𝑟1 )𝑙1 = 𝑟1 𝑙
𝐴 = 𝜋(𝑟1 𝑙 + 𝑟2 𝑙) or
where r = ½ (r1 + r2) is the average radius of the band. Now we apply this formula to our strategy. Consider
the surface shown in Figure 9.9, which is obtained by rotating the curve y = f (x), a x b, about the x-axis,
where f is positive and has a continuous derivative.
30
(a) Surface of revolution (b) Approximating band
Figure 9.9
In order to define its surface area, we divide the interval [a, b] into n subintervals with endpoints x0, x1, . . . ,
xn and equal width x, as we did in determining arc length. If yi = f (xi ), then the point Pi(xi, yi ) lies on the
curve.
The part of the surface between xi – 1 and xi is approximated by taking the line segment Pi – 1Pi and rotating it
about the x-axis. The result is a band with slant height l = | Pi – 1Pi | and average radius r = 1/2(yi – 1 + yi) so,
by Formula 2, its surface area is
𝑦𝑖−1 + 𝑦𝑖
2𝜋 |𝑃𝑖−1 𝑃𝑖 |
2
As in the proof, We have
2
|𝑃𝑖−1 𝑃𝑖 | = √1 + [𝑓 ′ (𝑥𝑖∗ )] ∆𝑥
Where 𝑥𝑖∗ is some number in [xi – 1, xi ]. When x is small, we have yi = f (xi) f (xi) and also yi – 1 = f (xi – 1) f
(xi), since f is continuous. Therefore
𝑦𝑖−1 + 𝑦𝑖 2
2𝜋 |𝑃𝑖−1 𝑃𝑖 | ≈ 2𝜋𝑓(𝑥𝑖∗ )√1 + [𝑓 ′ (𝑥𝑖∗ )] ∆𝑥
2
and so an approximation to what we think of as the area of the complete surface of revolution is
𝑛
2
∑ 2𝜋𝑓(𝑥𝑖∗ )√1 + [𝑓 ′ (𝑥𝑖∗ )] ∆𝑥 … (3)
𝑖=1
This approximation appears to become better as n ∞ and, recognizing (3) as a Riemann sum for the
function 𝑔(𝑥) = 2𝜋𝑓(𝑥)√1 + [𝑓 ′ (𝑥)]2 , we have
𝑛 𝑏
2
lim ∑ 2𝜋𝑓(𝑥𝑖∗ )√1 + [𝑓 ′ (𝑥𝑖∗ )] ∆𝑥 = ∫ 2𝜋𝑓(𝑥)√1 + [𝑓 ′ (𝑥)]2 𝑑𝑥
𝑛→∞ 𝑎
𝑖=1
31
Therefore, in the case where f is positive and has a continuous derivative, we define the surface area of the
surface obtained by rotating the curve y = f (x), a x b, about the x-axis as
If the curve is described as x = g(y), c y d, then the formula for surface area becomes
Now both Formulas 5 and 6 can be summarized symbolically, using the notation for arc length, as
For rotation about the y-axis, the surface area formula becomes
𝑑𝑦 2 𝑑𝑥 2
𝑑𝑠 = √1 + ( ) 𝑑𝑥 or 𝑑𝑠 = √1 + ( ) 𝑑𝑦
𝑑𝑥 𝑑𝑦
These formulas can be remembered by thinking of 2 y or 2 x as the circumference of a circle traced out
by the point (x, y) on the curve as it is rotated about the x-axis or y-axis, respectively (see Figure 9.10).
32
(a) Rotation about x-axis: S = ∫ 2 y ds (b) Rotation about y-axis: S = ∫ 2 x ds
Figure 9.10
Example
The curve 𝑦 = √4 − 𝑥 2 , − 1 ≤ 𝑥 ≤ 1, is an arc of the circle 𝑥 2 + 𝑦 2 = 4. Find the area of the surface
obtained by rotating this arc about the x-axis. (The surface is a portion of a sphere of radius 2. See Figure
9.11.)
Solution
We have
𝑑𝑦 1 1 −𝑥
= (4 − 𝑥 2 )−2 (−2𝑥) =
𝑑𝑥 2 √4 − 𝑥 2
1
𝑑𝑦 2
√
𝑆 = ∫ 2𝜋𝑦 1 + ( ) 𝑑𝑥
−1 𝑑𝑥
1
𝑥2
= 2𝜋 ∫ √4 − 𝑥 2 √1 + 𝑑𝑥
−1 4 − 𝑥2
33
1
4 − 𝑥2 + 𝑥2
= 2𝜋 ∫ √4 − 𝑥 2 √ 𝑑𝑥
−1 4 − 𝑥2
1 1
2
= 2𝜋 ∫ √4 − 𝑥 2 𝑑𝑥 = 4𝜋 ∫ 1 𝑑𝑥
−1 √4 − 𝑥 2 −1
= 4𝜋(2) = 8𝜋
Deep-sea divers realize that water pressure increases as they dive deeper. This is because the weight of the
water above them increases.
In general, suppose that a thin horizontal plate with area A square meters is submerged in a fluid of density
kilograms per cubic meter at a depth d meters below the surface of the fluid as in Figure 9.12.
The fluid directly above the plate has volume V = Ad, so its mass is m = V = Ad. The force exerted by the
fluid on the plate is therefore
𝐹 = 𝑚𝑔 = 𝜌𝑔𝐴𝑑
where g is the acceleration due to gravity. The pressure P on the plate is defined to be the force per unit
area:
𝐹
𝑃= = 𝜌𝑔𝑑
𝐴
The SI unit for measuring pressure is a newton per square meter, which is called a pascal
(abbreviation: 1 N/m2 = 1 Pa).
Since this is a small unit, the kilopascal (kPa) is often used. For instance, because the density of water is
= 1000 kg/m3, the pressure at the bottom of a swimming pool 2 m deep is
= 19,600 Pa
= 19.6 kPa
34
An important principle of fluid pressure is the experimentally verified fact that at any point in a liquid the
pressure is the same in all directions. (A diver feels the same pressure on nose and both ears.). Thus the
pressure in any direction at a depth d in a fluid with mass density is given by
𝑃 = 𝜌𝑔𝑑 = 𝛿𝑑 … (9)
This helps us determine the hydrostatic force against a vertical plate or wall or dam in a fluid. This is not a
straightforward problem because the pressure is not constant but increases as the depth increases.
Example
A dam has the shape of the trapezoid shown in Figure 9.13. The height is 20 m and the width is 50 m at the
top and 30 m at the bottom. Find the force on the dam due to hydrostatic pressure if the water level is 4 m
from the top of the dam.
We choose a vertical x-axis with origin at the surface of the water and directed downward as in Figure 3(a).
The depth of the water is 16 m, so we divide the interval [0, 16] into subintervals of equal length with
endpoints xi and we choose xi* [xi – 1, xi]. The ith horizontal strip of the dam is approximated by a
rectangle with height x and width wi, where, from similar triangles in Figure 9.14,
𝑎 10
∗ =
16 − 𝑥𝑖 20
Or
16 − 𝑥𝑖∗ 𝑥𝑖∗
𝑎= =8−
2 2
35
and so wi = 2(15 + a)
𝑥𝑖∗
𝑤𝑖 = 2(15 + 𝑎) = 2 (15 + 8 − )
2
= 46 − 𝑥𝑖∗
If x is small, then the pressure Pi on the i th strip is almost constant and we can use Eq. (9) to write
𝑃𝑖 ≈ 1000𝑔𝑥𝑖∗
The hydrostatic force Fi acting on the i th strip is the product of the pressure and the area:
Adding these forces and taking the limit as n ∞ , we obtain the total hydrostatic force on the dam:
𝑛 16
𝐹 = lim ∑ 1000𝑔𝑥𝑖∗ (46 − 𝑥𝑖∗ )∆𝑥 = ∫ 1000𝑔𝑥(46 − 𝑥)𝑑𝑥
𝑛→∞ 0
𝑖=1
16 16
2 )𝑑𝑥
𝑥3
= 1000(9.8) ∫ (46𝑥 − 𝑥 = 9000 [23𝑥 − ]| ≈ 4.43 𝑥 107 𝑁
2
0 3 0
Our main objective here is to find the point P on which a thin plate of any given shape balances horizontally
as in Figure 9.15. This point is called the center of mass (or center of gravity) of the plate.
We first consider the simpler situation illustrated in Figure 9.16, where two masses m1 and m2 are attached
to a rod of negligible mass on opposite sides of a fulcrum and at distances d1 and d2 from the fulcrum.
36
Figure 9.16: Simple illustration of the centre of mass of the plate
𝑚1 𝑑1 = 𝑚2 𝑑2 … . (10)
This is an experimental fact discovered by Archimedes and called the Law of the Lever. (Think of a lighter
person balancing a heavier one on a seesaw by sitting farther away from the center.)
Now suppose that the rod lies along the x-axis with m1 at x1 and m2 at x2 and the center of mass at 𝑥̅ . If we
compare Figures 9.16 and 9.17, we see that 𝑑1 = 𝑥̅ − 𝑥1 and 𝑑2 = 𝑥2 − 𝑥̅ and so Eq. (10) gives
𝑚1 (𝑥̅ − 𝑥1 ) = 𝑚2 (𝑥2 − 𝑥̅ )
𝑚1 𝑥̅ + 𝑚2 𝑥̅ = 𝑚1 𝑥1 + 𝑚2 𝑥2
𝑚1 𝑥1 + 𝑚2 𝑥2
𝑥̅ = … . (11)
𝑚1 + 𝑚2
The numbers m1x1 and m2x2 are called the moments of the masses m1 and m2 (with respect to the origin),
and Eq. (11) says that the center of mass 𝑥̅ is obtained by adding the moments of the masses and dividing
by the total mass m = m1 + m2.
In general, if we have a system of n particles with masses m1, m2, . . . , mn located at the points x1, x2, . . . , xn
on the x-axis, it can be shown similarly that the center of mass of the system is located at
∑𝑛
𝑖=1 𝑚𝑖 𝑥𝑖 ∑𝑛
𝑖=1 𝑚𝑖 𝑥𝑖
𝑥̅ = ∑𝑛
= … . (12)
𝑖=1 𝑚𝑖 𝑚
where m = mi is the total mass of the system, and the sum of the individual moments
37
𝑛
𝑀 = ∑ 𝑚𝑖 𝑥𝑖
𝑖=1
is called the moment of the system about the origin. Then Eq. (12) could be rewritten as 𝑚𝑥̅ = 𝑀, which
says that if the total mass were considered as being concentrated at the center of mass 𝑥̅ , then its moment
would be the same as the moment of the system. Now we consider a system of n particles with masses m1,
m2, . . . , mn located at the points (x1, y1), (x2, y2), . . . , (xn, yn) in the xy-plane as shown
in Figure 9.18.
By analogy with the one-dimensional case, we define the moment of the system about the y-axis to be
𝑛
𝑀𝑦 = ∑ 𝑚𝑖 𝑥𝑖 … (13)
𝑖=1
𝑀𝑥 = ∑ 𝑚𝑖 𝑦𝑖 … (14)
𝑖=1
Then My measures the tendency of the system to rotate about the y-axis and Mx measures the tendency to
rotate about the x-axis. As in the one-dimensional case, the coordinates (𝑥,
̅ 𝑦̅) of the center of mass are
given in terms of the moments by the formulas
𝑀𝑦 𝑀𝑥
𝑥̅ = 𝑦̅ = … (15)
𝑚 𝑚
where m = mi is the total mass. Since 𝑚𝑥̅ = 𝑀𝑦 and 𝑚𝑦̅ = 𝑀𝑥 , the center of mass (𝑥,
̅ 𝑦̅) is the point
where a single particle of mass m would have the same moments as the system.
Example
Find the moments and center of mass of the system of objects that have masses 3, 4, and 8 at the points (–
1, 1), (2, –1), and (3, 2), respectively.
38
Solution
𝑀𝑦 29
𝑥̅ = =
𝑚 15
𝑀𝑥 15
𝑦̅ = = =1
𝑚 15
14
Thus the center of mass is (1 , 1) (See Figure 9.19)
15
14
Figure 9.19: Centre of mass (1 15 , 1)
Next we consider a flat plate (called a lamina) with uniform density that occupies a region R of the plane.
We wish to locate the center of mass of the plate, which is called the centroid of R.
In doing so we use the following physical principles: The symmetry principle says that if R is symmetric
about a line l, then the centroid of R lies on l. (If R is reflected about l, then R remains the same so its
centroid remains fixed. But the only fixed points lie on l.). Thus the centroid of a rectangle is its center.
Moments should be defined so that if the entire mass of a region is concentrated at the center of mass,
then its moments remain unchanged. Also, the moment of the union of two nonoverlapping regions should
be the sum of the moments of the individual regions.
Suppose that the region R is of the type shown in Figure 9.20 (a) ; that is, R lies between the lines x = a and
x = b, above the x-axis, and beneath the graph of f, where f is a continuous function.
1
The centroid of the i th approximating rectangle Ri is its center 𝐶𝑖 (𝑥̅𝑖 , 2 𝑓(𝑥̅ )). Its area is 𝑓(𝑥̅𝑖 )∆𝑥, so its
mass is
𝜌𝑓(𝑥̅𝑖 )∆𝑥
The moment of Ri about the y-axis is the product of its mass and the distance from Ci to the y-axis, which is
𝑥̅𝑖 . Thus
Adding these moments, we obtain the moment of the polygonal approximation to R, and then by taking the
limit as n ∞, we obtain the moment of R itself about the y-axis:
In a similar fashion we compute the moment of Ri about the x-axis as the product of its mass and the
distance from Ci to the x-axis:
1 1
𝑀𝑥 (𝑅𝑖 ) = [𝜌𝑓(𝑥̅𝑖 )∆𝑥] 𝑓(𝑥̅𝑖 ) = 𝜌 ∙ [𝑓(𝑥̅𝑖 )]2 ∆𝑥
2 2
Again we add these moments and take the limit to obtain the moment of R about the x-axis:
40
Just as for systems of particles, the center of mass of the plate is defined so that 𝑚𝑥̅ = 𝑀𝑦 and 𝑚𝑦̅ = 𝑀𝑥 .
But the mass of the plate is the product of its density and its area:
𝑏
𝑚 = 𝜌𝐴 = 𝜌 ∫ 𝑓(𝑥)𝑑𝑥
𝑎
And so
Notice the cancellation of the ’s. The location of the center of mass is independent of the density. In
summary, the center of mass of the plate (or the centroid of R) is located at the point (𝑥̅ , 𝑦̅), where
… (16)
If the region R lies between two curves y = f (x) and y = g (x), where f (x) g (x), as illustrated in Figure 9.21,
then the same sort of argument that led to Formulas 16 can be used to show that the centroid of R is (𝑥̅ , 𝑦̅),
where
…. (17)
Example
A torus is formed by rotating a circle of radius r about a line in the plane of the circle that is a distance R (>
r) from the center of the circle. Find the volume of the torus.
Solution
The circle has area A = r2. By the symmetry principle, its centroid is its center and so the distance traveled
by the centroid during a rotation is d = 2 R.
𝑉 = 𝐴𝑑 = (2𝜋𝑅)(𝜋𝑟 2 ) = 2𝜋 2 𝑟 2 𝑅
42
MULTIPLE INTEGRALS
WEEK 10 : MULTIPLE INTEGRALS
10.1 DOUBLE INTEGRALS OVER RECTANGLES
10.1.1 REVIEW OF DEFINITE INTEGRALS
First let’s recall the basic facts concerning definite integrals of functions of a single variable. If f (x) is defined
for a x b, we start by dividing the interval [a, b] into n subintervals [xi – 1, xi] of equal width
x = (b – a)/n and we choose sample points 𝑥𝑖∗ in these subintervals. Then we form the Riemann sum
𝑛
and take the limit of such sums as n ∞ to obtain the definite integral of f from a to b:
𝑏 𝑛
In the special case where f (x) 0, the Riemann sum can be interpreted as the sum of the areas of the
𝑏
approximating rectangles in Figure 10.1, and ∫𝑎 𝑓(𝑥)𝑑𝑥 represents the area under the curve y = f (x) from a
to b.
and we first suppose that f (x, y) 0. The graph of f is a surface with equation z = f (x, y).
43
Let S be the solid that lies above R and under the graph of f, that is,
Our goal is to find the volume of S. The first step is to divide the rectangle R into subrectangles. We
accomplish this by dividing the interval [a, b] into m subintervals [xi – 1, xi] of equal width x = (b – a)/m and
dividing [c, d ] into n subintervals [yj – 1, yj] of equal width y = (d – c)/n.
By drawing lines parallel to the coordinate axes through the endpoints of these subintervals, as in Figure 3,
we form the subrectangles
If we choose a sample point (𝑥𝑖𝑗∗ , 𝑦𝑖𝑗∗ ) in each Rij, then we can approximate the part of S that lies above
each Rij by a thin rectangular box (or “column”) with base Rij and height (𝑥𝑖𝑗∗ , 𝑦𝑖𝑗∗ ) as shown in Figure 10.4.
44
The volume of this box is the height of the box times the area of the base rectangle:
If we follow this procedure for all the rectangles and add the volumes of the corresponding boxes, we get
an approximation to the total volume of S:
This double sum means that for each subrectangle we evaluate f at the chosen point and multiply by the
area of the subrectangle, and then we add the results. (See Figure 10.5.)
45
Our intuition tells us that the approximation given in Eq. (3) becomes better as m and n become larger and
so we would expect that
We use the expression in Eq. (4) to define the volume of the solid S that lies under the graph of f and above
the rectangle R. Limits of the type that appear in Equation 4 occur frequently, not just in finding volumes
but in a variety of other situations even when f is not a positive function. So we make the following
definition.
The precise meaning of the limit in Definition 5 is that for every number ε > 0 there is an integer N such that
𝑚 𝑛
for all integers m and n greater than N and for any choice of sample points (𝑥𝑖𝑗∗ , 𝑦𝑖𝑗∗ ) in Rij. A function f is
called integrable if the limit in Definition 5 exists.
It is shown in courses on advanced calculus that all continuous functions are integrable. In fact, the double
integral of f exists provided that f is “not too discontinuous.” In particular, if f is bounded on R, [that is,
there is a constant M such that | f (x, y) | M for all (x, y) in R], and f is continuous there, except on a finite
number of smooth curves, then f is integrable over R.
The sample point (𝑥𝑖𝑗∗ , 𝑦𝑖𝑗∗ ) can be chosen to be any point in the subrectangle Rij, but if we choose it to be
the upper right-hand corner of Rij [namely (xi, yj), see Figure 10.3], then the expression for the double
integral looks simpler:
By comparing Definitions 4 and 5, we see that a volume can be written as a double integral:
46
The sum in Definition 5,
𝑚 𝑛
∑ ∑ 𝑓(𝑥𝑖𝑗∗ , 𝑦𝑖𝑗∗ ) ∆𝐴
𝑖=1 𝑗=1
is called a double Riemann sum and is used as an approximation to the value of the double integral. [Notice
how similar it is to the Riemann sum in (1) for a function of a single variable.]. If f happens to be a positive
function, then the double Riemann sum represents the sum of volumes of columns, as in Figure 10.5, and is
an approximation to the volume under the graph of f.
Suppose that f is a function of two variables that is integrable on the rectangle R = [a, b] [c, d ].
𝑑
We use the notation ∫𝑐 𝑓(𝑥, 𝑦)𝑑𝑦 to mean that x is held fixed and f (x, y) is integrated with respect to y
from y = c to y = d. This procedure is called partial integration with respect to y. (Notice its similarity to
partial differentiation.)
𝑑
Now ∫𝑐 𝑓(𝑥, 𝑦)𝑑𝑦 is a number that depends on the value of x, so it defines a function of x:
𝑑
𝐴(𝑥) = ∫ 𝑓(𝑥, 𝑦)𝑑𝑦
𝑐
The integral on the right side of Eq. (7) is called an iterated integral. Usually the brackets are omitted. Thus
means that we first integrate with respect to y from c to d and then with respect to x from a to b. Similarly,
the iterated integral
means that we first integrate with respect to x (holding y fixed) from x = a to x = b and then we integrate
the resulting function of y with respect to y from y = c to y = d. Notice that in both Eq. (8) and (9) we work
from the inside out.
47
The following theorem gives a practical method for evaluating a double integral by expressing it as an
iterated integral (in either order).
In the special case where f (x, y) can be factored as the product of a function of x only and a function of y
only, the double integral of f can be written in a particularly simple form. To be specific, suppose that f (x, y)
= g (x)h (y) and R = [a, b] [c, d]. Then Fubini’s Theorem gives
𝑏
since ∫𝑎 𝑔(𝑥)𝑑𝑥 is a constant. Therefore, in this case, the double integral of f can be written as the product
of two single integrals:
48
Example
Find the volume of the solid that lies above the square R = [0, 2] [0, 2] and below the elliptic paraboloid z
= 16 – x2 – 2y2, as shown in Figure 10.6.
Solution
2 2
𝑥3
= ∫ [16𝑥 − − 2𝑥𝑦 2 ] 𝑑𝑦
0 3 0
2 2 2
88 88 4𝑦 3
= ∫ [ − 4𝑦 2 ] 𝑑𝑦 = [ 𝑦 − ] = 48 unit 3
0 3 0 3 3 0
2 2
2
= ∫ [16𝑦 − 𝑥 2 𝑦 − 𝑦 3 ] 𝑑𝑥
0 3 0
2 2 2
80 80 2𝑥 3
=∫ [ − 2𝑥 2 ] 𝑑𝑥 = [ 𝑥 − ] = 48 unit 3
0 3 0 3 3 0
49
Both method based on iterated integrals give the same answer.
Example
Solution
2 𝑦=2
2
𝑦2 2
∫ 𝑥 𝑦 𝑑𝑦 = [𝑥 ]
1 2 𝑦=1
22 12 3
= 𝑥2 ( ) − 𝑥2 ( ) = 𝑥2
2 2 2
3
Thus the function A in the preceding discussion is given by 𝐴(𝑥) = 𝑥 2 in this example. We now integrate
2
this function of x from 0 to 3:
3 2 3 2 3
3 2
∫ ∫ 𝑥 2 𝑦 𝑑𝑦 𝑑𝑥 = ∫ [∫ 𝑥 2 𝑦 𝑑𝑦] 𝑑𝑥 = ∫ 𝑥 𝑑𝑥
0 1 0 1 0 2
3
𝑥3 27
= | =
2 0 2
We know that the average value of a function f of one variable defined on an interval [a, b] is
𝑏
1
𝑓𝑎𝑣𝑒 = ∫ 𝑓(𝑥) 𝑑𝑥
𝑏−𝑎 𝑎
In a similar fashion we define the average value of a function f of two variables defined on a rectangle R to
be
1
𝑓𝑎𝑣𝑒 = ∬ 𝑓(𝑥, 𝑦) 𝑑𝐴
𝐴(𝑅)
𝑅
50
𝐴(𝑅) × 𝑓𝑎𝑣𝑒 = ∬ 𝑓(𝑥, 𝑦) 𝑑𝐴
𝑅
says that the box with base R and height fave has the same volume as the solid that lies under the graph of f.
[If z = f (x, y) describes a mountainous region and you chop off the tops of the mountains at height fave, then
you can use them to fill in the valleys so that the region becomes completely flat. See Figure 10.7.]
For single integrals, the region over which we integrate is always an interval. But for double integrals, we
want to be able to integrate a function f not just over rectangles but also over regions D of more general
shape, such as the one illustrated in Figure 10.8.
51
We suppose that D is a bounded region, which means that D can be enclosed in a rectangular region R as in
Figure 10.9.
… (12)
…. (13)
Where F is given by Eq. (12). Definition in Eq. (13) makes sense because R is a rectangle and so R F (x, y) dA
has been previously defined. The procedure that we have used is reasonable because the values of F (x, y)
are 0 when (x, y) lies outside D and so they contribute nothing to the integral.
This means that it doesn’t matter what rectangle R we use as long as it contains D. In the case where f (x, y)
0, we can still interpret D f (x, y) dA as the volume of the solid that lies above D and under the surface z =
f (x, y) (the graph of f ).
You can see that this is reasonable by comparing the graphs of f and F in Figures 10.10 and remembering
that R F (x, y) dA is the volume under the graph of F.
52
Figure 10.10 also shows that F is likely to have discontinuities at the boundary points of D. Nonetheless, if f
is continuous on D and the boundary curve of D is “well behaved”, then it can be shown that R F (x, y) dA
exists and therefore D f (x, y) dA exists. In particular, this is the case for type I and type II regions.
A plane region D is said to be of type I if it lies between the graphs of two continuous functions of x, that is,
In order to evaluate D f (x, y) dA when D is a region of type I, we choose a rectangle R = [a, b] [c, d ] that
contains D, as in Figure 10.12, and we let F be the function given by Eq. (12); that is, F agrees with f on D
and F is 0 outside D.
𝑏 𝑑
∬ 𝑓(𝑥, 𝑦) 𝑑𝐴 = ∬ 𝐹(𝑥, 𝑦) 𝑑𝐴 = ∫ ∫ 𝐹(𝑥, 𝑦)𝑑𝑦 𝑑𝑥
𝑎 𝑐
𝐷 𝑅
53
Observe that F (x, y) = 0 if y < g1 (x) or y > g2 (x) because (x, y) then lies outside D. Therefore
𝑑 𝑔2 (𝑥) 𝑔2 (𝑥)
∫ 𝐹(𝑥, 𝑦)𝑑𝑦 = ∫ 𝐹(𝑥, 𝑦)𝑑𝑦 = ∫ 𝑓(𝑥, 𝑦)𝑑𝑦
𝑐 𝑔1 (𝑥) 𝑔1 (𝑥)
because F (x, y) = f (x, y) when g1 (x) y g2 (x). Thus we have the following formula that enables us to
evaluate the double integral as an iterated integral.
.. (14)
The integral on the right side of Eq. (14) is an iterated integral, except that in the inner integral we regard x
as being constant not only in f (x, y) but also in the limits of integration, g1(x) and g2(x).
where h1 and h2 are continuous. Two such regions are illustrated in Figure 10.13.
Using the same methods that were used in establishing Eq. (14), we can show that
2 𝑑 ℎ (𝑦)
∬𝐷 𝑓(𝑥, 𝑦) 𝑑𝐴 = ∫𝑐 ∫ℎ (𝑦) 𝑓(𝑥, 𝑦)𝑑𝑥 𝑑𝑦 … (16)
1
54
Example
Evaluate D (x + 2y) dA, where D is the region bounded by the parabolas y = 2x2 and y = 1 + x2.
Solution
We note that the region D, sketched in Figure 10.14, is a type I region but not a type II region and we can
write D = {(x, y) | –1 x 1, 2x2 y 1 + x2}
Since the lower boundary is y = 2x2 and the upper boundary is y = 1 + x2, Eq. (14) gives
1 1+𝑥 2
∬(𝑥 + 2𝑦) 𝑑𝐴 = ∫ ∫ (𝑥 + 2𝑦)𝑑𝑦 𝑑𝑥
−1 2𝑥 2
𝐷
1 2
𝑦=1+𝑥
= ∫ [𝑥𝑦 + 𝑦 2 ]𝑦=2𝑥 2 𝑑𝑥
−1
1
= ∫ [𝑥(1 + 𝑥 2 ) + (1 + 𝑥 2 )2 − 𝑥(2𝑥 2 ) − (2𝑥 2 )2 ]𝑑𝑥
−1
1
32
= ∫(−3𝑥 4 − 𝑥 3 + 2𝑥 2 + 𝑥 + 1)𝑑𝑥 =
15
−1
55
10.2.1 PROPERTIES OF DOUBLE INTEGRALS
We assume that all of the following integrals exist. For rectangular regions D the first three properties can
be proved in the same manner. And then for general regions the properties follow from definition in Eq.
(13).
…(17)
…(18)
…(19)
The next property of double integrals is similar to the property of single integrals given by the equation
If D = D1 U D2, where D1 and D2 don’t overlap except perhaps on their boundaries (see Figure 10.15), then
…(20)
Property in Eq. (20) can be used to evaluate double integrals over regions D that are neither type I nor type
II but can be expressed as a union of regions of type I or type II.
56
Figure 10.16 illustrates this procedure.
Figure 10.16
The next property of integrals says that if we integrate the constant function f (x, y) = 1 over a region D, we
get the area of D:
Suppose that we want to evaluate a double integral R f (x, y) dA, where R is one of the regions shown in
Figure 10.17. In either case the description of R in terms of rectangular coordinates is rather complicated,
but R is easily described using polar coordinates.
Figure 10.17
57
From Figure 10.18, the polar coordinates (r, ) of a point are related to the rectangular coordinates (x, y) by
the equations
𝑟2 = 𝑥2 + 𝑦2 𝑥 = 𝑟 cos 𝜃 𝑦 = 𝑟 sin 𝜃
R = {(r, ) | a r b, }
In order to compute the double integral R f (x, y) dA, where R is a polar rectangle, we divide the interval [a,
b] into m subintervals [ri – 1, ri] of equal width r = (b – a)/m and we divide the interval [, ] into n
subintervals [j – 1, j] of equal width = ( – )/n. Then the circles r = ri and the rays = j divide the polar
rectangle R into the small polar rectangles Rij shown in Figure 10.20.
58
Figure 10.20: Dividing R into polar subrectangles
We compute the area of Rij using the fact that the area of a sector of a circle with radius r and central angle
1
is 2 𝑟 2 𝜃. Subtracting the areas of two such sectors, each of which has central angle = j – j – 1, we find
that the area of Rij is
1 1 2 1
∆𝐴𝑖 = 𝑟𝑖2 ∆𝜃 − 𝑟𝑖−1 ∆𝜃 = (𝑟𝑖2 − 𝑟𝑖−1
2
)∆𝜃
2 2 2
1
= (𝑟𝑖 + 𝑟𝑖−1 )(𝑟𝑖 − 𝑟𝑖−1 )∆𝜃 = 𝑟𝑖∗ ∆𝑟 ∆𝜃
2
Although we have defined the double integral R f (x, y) dA in terms of ordinary rectangles, it can be shown
that, for continuous functions f, we always obtain the same answer using polar rectangles. The rectangular
coordinates of the center of Rij are (𝑟𝑖∗ cos 𝜃𝑗∗ , 𝑟𝑖∗ sin 𝜃𝑗∗ , ), so a typical Riemann sum is
… (21)
If we write g(r, ) = r f (r cos , r sin ), then the Riemann sum in Eq. (21) can be written as
𝑚 𝑛
59
𝛽 𝑏
∫ ∫ 𝑔(𝑟, 𝜃) 𝑑𝑟 𝑑𝜃
𝛼 𝑎
Therefore we have
𝑚 𝑛
∬ 𝑓(𝑥, 𝑦)𝑑𝐴 = lim ∑ ∑ 𝑓(𝑟𝑖∗ cos 𝜃𝑗∗ , 𝑟𝑖∗ sin 𝜃𝑗∗ ) ∆𝐴𝑖
𝑚,𝑛→∞
𝑅 𝑖=1 𝑗=1
𝑚 𝑛 𝛽 𝑏
= lim ∑ ∑ 𝑔(𝑟𝑖∗ , 𝜃𝑗∗ )∆𝑟 ∆𝜃 = ∫ ∫ 𝑔(𝑟, 𝜃) 𝑑𝑟 𝑑𝜃
𝑚,𝑛→∞ 𝛼 𝑎
𝑖=1 𝑗=1
𝛽 𝑏
= ∫ ∫ 𝑓(𝑟 cos 𝜃 , 𝑟 sin 𝜃) 𝑟 𝑑𝑟 𝑑𝜃
𝛼 𝑎
…(22)
The formula in Eq. (22) says that we convert from rectangular to polar coordinates in a double integral by
writing x = r cos and y = r sin , using the appropriate limits of integration for r and , and replacing dA by
r dr d.
Be careful not to forget the additional factor r on the right side of Formula in Eq. (22). A classical method for
remembering this is shown in Figure 10.21, where the “infinitesimal” polar rectangle can be thought of as
an ordinary rectangle with dimensions r d and dr and therefore has “area” dA = r dr d.
60
Example
Evaluate R (3x + 4y2) dA, where R is the region in the upper half-plane bounded by the circles x2 + y2 = 1
and x2 + y2 = 4.
Solution
R = {(x, y) | y 0, 1 x2 + y2 4}
𝜋 2
2 )𝑑𝐴
∬(3𝑥 + 4𝑦 = ∫ ∫ (3𝑟 cos 𝜃 + 4𝑟 2 sin2 𝜃) 𝑟 𝑑𝑟 𝑑𝜃
0 1
𝑅
𝜋 2 𝜋
= ∫ ∫ (3𝑟 2 cos 𝜃 + 4𝑟 3 sin2 𝜃) 𝑑𝑟 𝑑𝜃 = ∫ [𝑟 3 cos 𝜃 + 𝑟 4 sin2 𝜃]𝑟=2
𝑟=1 𝑑𝜃
0 1 0
𝜋 𝜋
15
= ∫ (7 cos 𝜃 + 15sin2 𝜃 ) 𝑑𝜃 = ∫ [7 cos 𝜃 + (1 − cos 2𝜃) ] 𝑑𝜃
0 0 2
𝜋
15𝜃 15 15𝜋
= 7 sin 𝜃 + − sin 2𝜃] =
2 4 0 2
What we have done so far can be extended to the more complicated type of region shown in Figure 10.23.
In fact, by combining formula in Eq. (22) with
𝑑 ℎ2 (𝑦)
∬ 𝑓(𝑥, 𝑦)𝑑𝐴 = ∫ ∫ 𝑓(𝑥, 𝑦)𝑑𝑥 𝑑𝑦
𝑐 ℎ1 (𝑦)
𝐷
61
… (23)
In particular, taking f (x, y) = 1, h1( ) = 0, and h2( ) = h( ) in this formula, we see that the area of the region
D bounded by = , = , and r = h( ) is
𝛽 ℎ(𝜃)
𝐴𝐷 = ∬ 1 𝑑𝐴 = ∫ ∫ 𝑟 𝑑𝑟 𝑑𝜃
𝛼 0
𝐷
𝛽 ℎ(𝜃)
𝑟2 𝛽
1
=∫ [ ] 𝑑𝜃 = ∫ [ℎ(𝜃)]2 𝑑𝜃
𝛼 2 0 𝛼 2
We have defined single integrals for functions of one variable and double integrals for functions of two
variables, so we can define triple integrals for functions of three variables. Let’s first deal with the simplest
case where f is defined on a rectangular box:
B = {(x, y, z) | a x b, c y d, r z s} … (24)
The first step is to divide B into sub-boxes. We do this by dividing the interval [a, b] into l subintervals [xi–1,
xi] of equal width x, dividing [c, d ] into m subintervals of width y, and dividing [r, s] into n subintervals of
width z.
The planes through the endpoints of these subintervals parallel to the coordinate planes divide the box B
into lmn sub-boxes
∗ ∗ ∗
where the sample point (𝑥𝑖𝑗𝑘 , 𝑦𝑖𝑗𝑘 , 𝑧𝑖𝑗𝑘 ) is in Bijk. By analogy with the definition of a double integral, we
define the triple integral as the limit of the triple Riemann sums in Eq. (25).
…(26)
Again, the triple integral always exists if f is continuous. We can choose the sample point to be any point in
the sub-box, but if we choose it to be the point (xi, yj, zk) we get a simpler-looking expression for the triple
integral:
63
Just as for double integrals, the practical method for evaluating triple integrals is to express them as
iterated integrals as follows.
…(27)
The iterated integral on the right side of Fubini’s Theorem means that we integrate first with respect to x
(keeping y and z fixed), then we integrate with respect to y (keeping z fixed), and finally we integrate with
respect to z.
There are five other possible orders in which we can integrate, all of which give the same value. For
instance, if we integrate with respect to y, then z, and then x, we have
𝑏 𝑠 𝑑
∭ 𝑓(𝑥, 𝑦, 𝑧) 𝑑𝑉 = ∫ ∫ ∫ 𝑓(𝑥, 𝑦, 𝑧)𝑑𝑦 𝑑𝑧 𝑑𝑥
𝑎 𝑟 𝑐
𝐵
Example
Evaluate the triple integral B xyz2 dV, where B is the rectangular box given by
B = {(x, y, z) | 0 x 1, –1 y 2, 0 z 3}
Solution
We could use any of the six possible orders of integration. If we choose to integrate with respect to x, then
y, and then z, we obtain
3 2 1 3 2 𝑥=1
2 2
𝑥 2 𝑦𝑧 2
∭ 𝑥𝑦𝑧 𝑑𝑉 = ∫ ∫ ∫ 𝑥𝑦𝑧 𝑑𝑥 𝑑𝑦 𝑑𝑧 = ∫ ∫ [ ] 𝑑𝑦 𝑑𝑧
0 −1 0 0 −1 2 𝑥=0
𝐵
3 2 𝑦=2
𝑦𝑧 2 3
𝑦2𝑧2 3
3𝑧 2
=∫ ∫ 𝑑𝑦 𝑑𝑧 = ∫ [ ] 𝑑𝑧 = ∫ 𝑑𝑧
0 −1 2 0 4 𝑦=−1 0 4
3
𝑧3 27
= ] =
4 0 4
64
10.4.1 TRIPLE INTEGRALS OVER A GENERAL BOUNDED REGION E
Now we define the triple integral over a general bounded region E in three-dimensional space (a solid) by
much the same procedure that we used for double integrals. We enclose E in a box B of the type given by
Eq. (24). Then we define F so that it agrees with f on E but is 0 for points in B that are outside E. By
definition
This integral exists if f is continuous and the boundary of E is “reasonably smooth”. The triple integral has
essentially the same properties as the double integral. We restrict our attention to continuous functions f
and to certain simple types of regions.
A solid region E is said to be of type 1 if it lies between the graphs of two continuous functions of x and y,
that is,
Notice that the upper boundary of the solid E is the surface with equation z = u2(x, y), while the lower
boundary is the surface z = u1(x, y). By the same sort of argument, it can be shown that if E is a type 1 region
given by Eq. (28), then
…(29)
The meaning of the inner integral on the right side of Eq. (29) is that x and y are held fixed, and therefore
u1(x, y) and u2(x, y) are regarded as constants, while f (x, y, z) is integrated with respect to z.
65
In particular, if the projection D of E onto the xy-plane is a type I plane region (as in Figure 10.26).
Figure 10.26: A type I solid region where the projection D is a type I plane region
Then
…(30)
If, on the other hand, D is a type II plane region (as in Figure 10.27), then
…(31)
where, this time, D is the projection of E onto the yz-plane (see Figure 10.28).
The back surface is x = u1(y, z), the front surface is x = u2(y, z), and we have
..(32)
where D is the projection of E onto the xz-plane, y = u1(x, z) is the left surface, and y = u2(x, z) is the right
surface (see Figure 10.29).
67
For this type of region we have
…(33)
In each of Eq. (32) and (33) there may be two possible expressions for the integral depending on whether D
is a type I or type II plane region (and corresponding to Eq. (30) and (31)).
In plane geometry the polar coordinate system is used to give a convenient description of certain curves
and regions. Figure 10.30 enables us to recall the connection between polar and Cartesian coordinates.
If the point P has Cartesian coordinates (x, y) and polar coordinates (r, ), then, from the figure,
x = r cos y = r sin
𝑦
r2 = x2 + y2 tan = 𝑥
In three dimensions there is a coordinate system, called cylindrical coordinates, that is similar to polar
coordinates and gives convenient descriptions of some commonly occurring surfaces and solids. As we will
see, some triple integrals are much easier to evaluate in cylindrical coordinates.
In the cylindrical coordinate system, a point P in three-dimensional space is represented by the ordered
triple (r, , z), where r and are polar coordinates of the projection of P onto the xy-plane and z is the
directed distance from the xy-plane to P. (See Figure 10.31)
68
Figure 10.31: The cylindrical coordinates of a point
…(34)
…(35)
Example
Plot the point with cylindrical coordinates (2, 2 /3, 1) and find its rectangular coordinates.
Solution
The point with cylindrical coordinates (2, 2 /3, 1) is plotted in Figure 10.32. From Eq. (34), its rectangular
coordinates are
2𝜋 1
𝑥 = 2 cos = 2 (− ) = −1
3 2
2𝜋 √3
𝑦 = 2 sin = 2 ( ) = √3
3 2
𝑧=1
69
Figure 10.32: Cylindrical coordinates (2, 2 /3, 1)
Suppose that E is a type 1 region whose projection D onto the xy-plane is conveniently described in polar
coordinates (see Figure 10.33).
We know
… (36)
..(37)
70
Formula in Eq. (37) is the formula for triple integration in cylindrical coordinates.
It says that we convert a triple integral from rectangular to cylindrical coordinates by writing x = r cos , y =
r sin , leaving z as it is, using the appropriate limits of integration for z, r, and , and replacing dV by r dz dr
d. (Figure 10.34 shows how to remember this.)
It is worthwhile to use this formula when E is a solid region easily described in cylindrical coordinates, and
especially when the function f (x, y, z) involves the expression x2 + y2.
Example
A solid E lies within the cylinder x2 + y2 = 1, below the plane z = 4, and above the paraboloid z = 1 – x2 – y2.
(See Figure 10.35.) The density at any point is proportional to its distance from the axis of the cylinder. Find
the mass of E.
Solution
In cylindrical coordinates the cylinder is r = 1 and the paraboloid is z = 1 – r2, so we can write
E = {(r, , z) | 0 2, 0 r 1, 1 – r2 z 4}
Since the density at (x, y, z) is proportional to the distance from the z-axis, the density function is
𝑓(𝑥, 𝑦, 𝑧) = 𝐾√𝑥 2 + 𝑦 2 = 𝐾𝑟
71
Therefore, the mass of E is
2𝜋 1 4
𝑚 = ∭ 𝐾√𝑥 2 + 𝑦 2 𝑑𝑉 = ∫ ∫ ∫ (𝐾𝑟)𝑟 𝑑𝑧 𝑑𝑟 𝑑𝜃
0 0 1−𝑟 2
𝐸
2𝜋 1 2𝜋 1
2 [4 2
=∫ ∫ 𝐾𝑟 − (1 − 𝑟 )] 𝑑𝑟 𝑑𝜃 = 𝐾 ∫ 𝑑𝜃 ∫ (3𝑟 2 + 𝑟 4 )𝑑𝑟
0 0 0 0
𝑟5 1 12𝜋𝐾
= 2𝜋𝐾 [𝑟 3 + ] =
5 0 5
72
DIFFERENTIAL EQUATIONS
WEEK 11: DIFFERENTIAL EQUATIONS
11.1 INTRODUCTION TO DIFFERENTIAL EQUATION
In creating a mathematical model of a physical system, we frequently involve differential equation/ integral
equation / integro-differential equations to express relationships, such as ‘the force acting on a falling object
is proportional to its acceleration’, ‘voltage drop across a resistor is proportional to the current’, etc.
Table 12.2: Three types of equations of a mathematical model: (i) Differential equation (ii) Integral
equation (iii) Integro-differential equation.
(i) Differential equation (ii) Integral equation (iii) Integro-differential equation
Equations which involve Equations which involve Equations which involve both
derivatives of the variables in integrals of the variables in the derivatives and integrals of the
the model. model. variables in the model.
From experiments, we know that the voltage loss through a resistor and capacitor is proportional to the
current and charge respectively, where Δ𝑉𝑟𝑒𝑠𝑖𝑠𝑡𝑜𝑟 (𝑡) ∝ 𝑖(𝑡) and Δ𝑉𝑐𝑎𝑝𝑎𝑐𝑖𝑡𝑜𝑟 (𝑡) ∝ 𝑞(𝑡). Hence, Δ𝑉𝑅 (𝑡) =
1
𝑅𝑖(𝑡) and Δ𝑉𝐶 (𝑡) = 𝑞(𝑡).
𝐶
According to Kirchhoff’s voltage law, the summation of voltage in a closed loop is equal to zero. Thus,
(𝑉𝐵 − 𝑉𝐴 ) + (𝑉𝐹 − 𝑉𝐵 ) + (𝑉𝐷 − 𝑉𝐹 ) + (𝑉𝐴 − 𝑉𝐷 ) = 0
1
𝑅𝑖(𝑡) + 𝑞(𝑡) + 0 + (−𝐸(𝑡)) = 0
𝐶
1
𝑅𝑖(𝑡) + 𝑞(𝑡) = 𝐸(𝑡)
𝐶
73
From definition, the current is equal to the rate of charge flow or the charge is the integral of the current
over time.
𝑑𝑞(𝑡)
𝑖(𝑡) = 𝑑𝑡
or 𝑞(𝑡) = ∫ 𝑖(𝑡)𝑑𝑡
Rearrange it, we obtain three different forms of mathematical model as shown below, to solve the desired
variables (i.e. charge and current). Note that different methods and strategies are used to solve these
equations. In this study, we will focus on the topic of differential equation.
𝑑𝑞(𝑡) 1 1 𝑑𝑞(𝑡) 1
(i) 𝑅 + 𝐶 𝑞(𝑡) = 𝐸(𝑡) (ii) 𝑅𝑖(𝑡) + 𝐶 ∫ 𝑖(𝑡)𝑑𝑡 = 𝐸(𝑡) (iii) 𝑅 + 𝐶 ∫ 𝑖(𝑡)𝑑𝑡 = 𝐸(𝑡)
𝑑𝑡 𝑑𝑡
-involves derivative,
𝑑𝑞(𝑡)
. -involves integrals, ∫ 𝑖(𝑡)𝑑𝑡. -involves both derivative,
𝑑𝑞(𝑡)
and
𝑑𝑡 𝑑𝑡
-This is known as integral integrals, ∫ 𝑖(𝑡)𝑑𝑡.
-This is known as differential
equation.
equation. -This is known as integro-differential
equation.
Differential equation (DE) plays a fundamental role in engineering because many physical
phenomena are best formulated mathematically in terms of their rate of change. By solving a differential
equation, we can gain a deeper understanding of the physical processes that these equations are describing.
Some examples of fundamental laws that are written in terms of the rate of change of variables are shown
in the table below.
Different types of differential equations may require different strategies to solve the problem. Thus, it is
important for the user to understand, recognize and classify the correct categories of differential equations.
A differential equation expresses such that the dependent variable(s) depends on the independent variable.
Example (1):
𝑑2 𝑦 𝑑𝑦
𝑑𝑥 2
− 4𝑥 𝑑𝑥 = 𝑐𝑜𝑠2𝑥
This differential equation has dependent variable of 𝑦 and independent variable of 𝑥.
Note that the variable 𝑦 is in the function of 𝑥, i.e. 𝑦 = 𝑦(𝑥). In other words, 𝑦 is changed with respect to
𝑥.
Example (2):
𝑑2 𝑥 𝑑𝑥
𝑑𝑡 2
− 4𝑥 𝑑𝑡 = 𝑐𝑜𝑠2𝑡
This differential equation has dependent variable of 𝑥 and independent variable of 𝑡.
Note that the variable 𝑥 is in the function of 𝑡, i.e. 𝑥 = 𝑥(𝑡). In other words, 𝑥 is changed with respect to
𝑡.
75
(ii) Ordinary Differential Equation (ODE) versus Partial Differential Equation (PDE)
Differential equation can be categorized into 2 cases: ODE & PDE. The classification of ODE and PDE
depends on the number of independent variable, regardless of the number of dependent variables.
CASE 1: ODE
Those equations that involve ordinary derivatives (i.e. 𝑑 symbol) are called ODE. ODE has only one
independent variable. It can be separated into the ODE problem or system of ODE problem depends on the
number of dependent variable.
(i) One dependent variable (ii) More than one dependent variable
For example: Brine mixture problem For example: Population of rabbit & fox
𝑑𝑥 𝑥 𝑑𝑥
=2− , = 𝑎𝑥 − 𝑏𝑥𝑦,
𝑑𝑡 5 𝑑𝑡
CASE 2: PDE
Those equations that involve partial derivatives (i.e. ∂ symbol) are called PDE. PDE has two independent
variables or more. It can be separated into the PDE problem or system of PDE problem depends on the
number of dependent variable.
(i) One dependent variable (ii) More than one dependent variable
For example: Transient heat equation For example: Incompressible Navier-Stokes
∂𝑇(𝑥,𝑡) ∂2 𝑇(𝑥,𝑡)
equation for pipe flow
−𝛼 ∂𝑥 2
= 0,
∂𝑡 ∂𝑢𝑖 ∂𝑢 ∂2 𝑢𝑖 ∂ω
∂𝑡
+ 𝑢𝑗 ∂𝑥 𝑖 − 𝑣 ∂𝑥 + ∂𝑥 = 𝑔𝑖 ,
where 𝑇 = temperature; 𝛼 = thermal diffusivity. 𝑗 𝑗 ∂𝑥𝑗 𝑖
(i) More than one independent variables (x & t) (i) More than one independent variables (t , 𝑥𝑗 & 𝑥𝑗 )
(ii) One dependent variable (T) (ii) More than one dependent variables (𝑢𝑖 & ω)
Comment: PDE is out of scope in this study. It will be covered in KIX1002. At current stage, students
should know how to classify between ODE and PDE.
76
(iii) Order of a differential equation
The order of a differential equation is the degree of the highest derivative that occurs in the equation. The
approach to find the order of ODE and PDE is the same as illustrated below.
Example (2):
𝑑𝑥 𝑑𝑦
( )2 − 3 − 𝑥 + 𝑦 = cos(2𝑡)
𝑑𝑡 𝑑𝑡
Example (2):
𝜕𝑥 𝜕𝑥
𝜕𝑡
+ (𝜕𝑦)4 = cos(2𝑡) + 2𝑦
Example (2):
∂2 𝑓 𝜕𝑓
+ ( )3 = cos(2𝑡) + 2𝑦
∂t ∂y 𝜕𝑦
77
Note: The order of an equation is not affected by any power to which the derivatives may be raised.
Moreover, degree of a differential equation is the power of the highest order derivative.
Example (1):
𝑑𝑥 𝑑𝑦
( )2 − 3 − 𝑥 + 𝑦 = cos(2𝑡)
𝑑𝑡 𝑑𝑡
The differential equation above is a 1st order ODE with degree 2.
Example (2):
∂2 𝑓 𝜕𝑓
+ ( )3 = cos(2𝑡) + 2𝑦
∂t ∂y 𝜕𝑦
(ii) PDE because it has more than one independent variables (𝑡 & 𝑦)
∂2 𝑓
(iii) Degree 1 because the power of the highest order derivative ∂t ∂y is 1
Example (3):
𝜕𝑥 𝜕𝑥
+ ( )4 = cos(2𝑡) + 2𝑦
𝜕𝑡 𝜕𝑦
(ii) PDE because it has more than one independent variables (𝑡 & 𝑦)
𝜕𝑥
(iii) Degree 4 because the power of the highest order derivative (𝜕𝑦)4 is 4
78
(iv) Linear and nonlinear differential equations
We used to plot a linear graph of 𝑦1 versus 𝑦2 using 𝑦1 = 𝑚𝑦2 + 𝑐 (Linear Algebraic Eqn), where 𝑚 & 𝑐 are
the slope and the intercept respectively. In this case, 𝑦1 is in the function of 𝑦2 . Therefore, 𝑦1 is the
dependent variable while 𝑦2 is the independent variable. Rearrange it, we obtain the general form of linear
equation as follows: 𝑎1 𝑦1 + 𝑎0 𝑦2 = 𝑐 where 𝑎1 = 1 and 𝑎0 = −𝑚.
By using similar approach, we get 1st order linear ODE where 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥). In this case,
𝑦′ is the first derivative of 𝑦. Linear ODE has the properties of 𝑓(𝑦1 + 𝑦2 ) = 𝑓(𝑦1 ) + 𝑓(𝑦2 ).
Any equation of ODE that does not follow the linear format as equation above is known as nonlinear ODE.
For example:
𝑎𝑛 (𝑥, 𝑦)(𝑦 (𝑛) ) 𝐴 + 𝑎𝑛−1 (𝑥, 𝑦)(𝑦 (𝑛−1) )𝐵 + ⋯ + 𝑎1 (𝑥, 𝑦)(𝑦 ′ )𝐶 + 𝑎0 (𝑥, 𝑦)(𝑦)𝐷 = 𝑔(𝑥, 𝑦).
where the power 𝐴, 𝐵, 𝐶 & 𝐷 ≠ 1
79
General format of 2nd order linear ODE: 𝑎2 (𝑥)𝑦 ′′ + 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥)
For example:
𝑑2 𝑓 𝑑𝑓
2nd order linear ODE: 𝑑𝑥 2 − 4𝑥 𝑑𝑥 − 𝑐𝑜𝑠2𝑥 − 3 = 0,
𝑑2 𝑓 𝑑𝑓
Rearrange it into the general format: 𝑑𝑥2 − 4𝑥 𝑑𝑥 = 𝑐𝑜𝑠2𝑥 + 3,
where
𝑑2 𝑓
𝑓′′ = 𝑑𝑥 2 ; (𝑓 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒 & 𝑥 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒)
𝑑𝑓
𝑓′ = ;
𝑑𝑥
𝑎2 (𝑥) = 1;
𝑎1 (𝑥) = −4𝑥;
𝑎0 (𝑥) = 0;
𝑔(𝑥) = 𝑐𝑜𝑠2𝑥 + 3
∴ It is a linear ODE since it follows the linear format: 𝑎2 (𝑥)𝑓 ′′ + 𝑎1 (𝑥)𝑓 ′ + 𝑎0 (𝑥)𝑓 = 𝑔(𝑥)
General format of 3rd order linear ODE: 𝑎3 (𝑥)𝑦 ′′′ + 𝑎2 (𝑥)𝑦 ′′ + 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥)
For example:
𝑑3𝑥 𝑑𝑥
3rd order linear ODE: 4 +3 + 2𝑥 + cos(𝑡) = 2,
𝑑𝑡 3 𝑑𝑡
𝑑3 𝑥 𝑑𝑥
Rearrange it into the general format: 4 𝑑𝑡 3 + 3 𝑑𝑡 + 2𝑥 = 2 − cos(𝑡),
where
𝑑3 𝑥
𝑥′′′ = 𝑑𝑡 3
; (𝑥 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒 & 𝑡 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒)
𝑑𝑥
𝑥′ = 𝑑𝑡
;
𝑎3 (𝑡) = 4;
𝑎2 (𝑡) = 0;
𝑎1 (𝑡) = 3;
𝑎0 (𝑡) = 2;
𝑔(𝑡) = 2 − cos(𝑡)
∴ It is a linear ODE since it follows the linear format: 𝑎3 (𝑡)𝑥 ′′′ + 𝑎2 (𝑡)𝑥 ′′ + 𝑎1 (𝑡)𝑥 ′ + 𝑎0 (𝑡)𝑥 = 𝑔(𝑡)
80
Case 2: Nonlinear ODE
General format of 1st order linear ODE: 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥)
For example:
𝑑𝑦 2
1st order nonlinear ODE: −4 (𝑑𝑥 ) = 𝑥 2 ,
∴ It is a nonlinear ODE because it does not obey linear equation: 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥) as the
𝑑𝑦 2 𝑑𝑦
derivative 𝑦′ ≠ ( ) where the is squared.
𝑑𝑥 𝑑𝑥
General format of 2nd order linear ODE: 𝑎2 (𝑥)𝑦 ′′ + 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥)
For example:
𝑑2 𝑓 𝑑𝑓
2nd order nonlinear ODE: − 4𝑓 + 𝑐𝑜𝑠2𝑥 = 0,
𝑑𝑥 2 𝑑𝑥
𝑑2 𝑓 𝑑𝑓
Rearrange it into the general format: 𝑑𝑥2 − 4𝑓 𝑑𝑥 = −𝑐𝑜𝑠2𝑥
∴ It is a nonlinear ODE because it does not obey linear equation: 𝑎2 (𝑥)𝑓 ′′ + 𝑎1 (𝑥)𝑓 ′ + 𝑎0 (𝑥)𝑓 = 𝑔(𝑥) as
the 𝑎1 (𝑥) ≠ −4𝑓, where 𝑎1 (𝑥) should not be in the function of dependent variable 𝑓.
General format of 3rd order linear ODE: 𝑎3 (𝑥)𝑦 ′′′ + 𝑎2 (𝑥)𝑦 ′′ + 𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥)
For example:
𝑑3 𝑥 𝑑𝑥
3rd order nonlinear ODE: 4 𝑑𝑡 3 + 3 𝑑𝑡 + 2sin(𝑥) + cos(𝑡) = 2,
𝑑3 𝑥 𝑑𝑥
Rearrange it into the general format: 4 𝑑𝑡 3 + 3 𝑑𝑡 + 2 sin(𝑥) = 2 − cos(𝑡),
∴ It is a nonlinear ODE because it does not obey linear equation: 𝑎3 (𝑡)𝑥 ′′′ + 𝑎2 (𝑡)𝑥 ′′ + 𝑎1 (𝑡)𝑥 ′ +
𝑎0 (𝑡)𝑥 = 𝑔(𝑡) as the 𝑥 ≠ sin(𝑥), where it has nonlinear sinusoidal function of dependent function 𝑥.
Hint: Most of the time, nonlinear ODE has nonlinear components such as coefficient 𝑎𝑖 (𝑥, 𝑦) in the function
𝑑𝑦
of dependent variable 𝑦 and the 𝑦 or its derivative have degree more than one, i.e. 𝑦 2 & (𝑑𝑥 )3 . For linear
differential equations, there are no products of the dependent variable and its derivatives and neither the
derivative occur to any power other than the first power.
81
For your additional knowledge, many of the nonlinear equations that occur in engineering cannot
be solved easily as they stand, but can be solved, for practical engineering purpose, by the process of
replacing them with linear equations that are a close approximation – at least in some region of interest.
This is applied to the case of linear differential equation only. Arrange the linear equation in standard format,
where all terms containing dependent variable occur on left-hand side (LHS), while terms containing only
the independent variable and constant occur on the right-hand side (RHS). Linear ODE can be categorized
into homogenous and non-homogeneous equation by evaluating the RHS term as follows.
Example (1):
𝑑𝑥
𝑑𝑡
+ 4𝑥 = 0
Example (2):
𝑑2 𝑥 𝑑𝑥
+3 + 𝑥𝑠𝑖𝑛(𝑡) = 0
𝑑𝑡 2 𝑑𝑡
Example (1):
𝜕𝑓 𝜕𝑓
𝜕𝑥
+ 𝜕𝑦 = 0
Example (2):
𝜕2 𝑓 𝜕𝑓
𝜕𝑥𝜕𝑦
+ 𝜕𝑦 = 0
82
Case 2: Nonhomogeneous equation
RHS term is equal to nonzero in the standard format
(i) Linear nonhomogeneous ODE
Example (1):
𝑑𝑥
𝑑𝑡
+ 4𝑥 = 5
Example (2):
𝑑2 𝑥 𝑑𝑥
𝑑𝑡 2
+ 3 𝑑𝑡 + 𝑥𝑠𝑖𝑛(𝑡) = cos(𝑡)
Example (1):
𝜕𝑓 𝜕𝑓
𝜕𝑥
+ 𝜕𝑦 = 4𝑥 2 + 2𝑦
Example (2):
𝜕2 𝑓 𝜕𝑓
+ = 5𝑦
𝜕𝑥𝜕𝑦 𝜕𝑦
Note: Other method is used to classify the homogeneity of nonlinear ODE as shown in section (vi).
𝑑𝑦 𝑓(𝑥,𝑦)
(vi) Homogeneous and nonhomogeneous equations of (𝑑𝑥 = 𝑔(𝑥,𝑦))
This section is particular important especially when we deal with 1st order nonlinear ODE problem. It is
worthwhile to mention that there is another method to classify the homogeneous and nonhomogeneous
groups in ODE. For 1st order ODE equation, the classification is given below.
First of all, the descriptions of the homogeneous and nonhomogeneous functions are given:
A function 𝑓(𝑥, 𝑦) is said to be homogeneous of degree 𝑛 if we get 𝑓(𝜆𝑥, 𝜆𝑦) = 𝜆𝑛 𝑓(𝑥, 𝑦) for all arbitrary
constant 𝜆.
Any function that does not follow the homogeneous format as equation above is known as a
nonhomogeneous function.
83
Example (1): Check the homogeneous degree for the function 𝑥 4 + 𝑥𝑦 3 .
Since 𝑓(𝜆𝑥, 𝜆𝑦) = 𝜆𝑛 𝑓(𝑥, 𝑦), thus we say the function 𝑥 4 + 𝑥𝑦 3 is homogeneous of degree, 𝑛 = 4.
Since 𝑓(𝜆𝑥, 𝜆𝑦) ≠ 𝜆𝑛 𝑓(𝑥, 𝑦) or (𝜆𝑦)2 − (𝜆𝑥)(𝜆𝑦) + 1 ≠ 𝜆2 (𝑦 2 − 𝑥𝑦 + 1) , thus we say the function
(𝑦)2 − (𝑥𝑦) + 1 is nonhomogeneous.
Then, the homogeneous and nonhomogeneous differential equation of 1st order ODE are given:
𝒅𝒚 𝒇(𝒙,𝒚)
Case (1): Homogeneous differential equation of =
𝒅𝒙 𝒈(𝒙,𝒚)
𝑑𝑦 𝑓(𝑥,𝑦)
𝑑𝑥
= 𝑔(𝑥,𝑦) is a homogeneous differential equation if 𝑓(𝑥, 𝑦) and 𝑔(𝑥, 𝑦) are homogeneous of the same
degree.
For example:
𝑑𝑦 𝑦 2 −𝑥 2
𝑑𝑥
= 2𝑥𝑦
It is a 1st order nonlinear homogeneous differential equation where the functions at numerator and
denominator are homogeneous of degree 2.
𝒅𝒚 𝒇(𝒙,𝒚)
Case (2): Nonhomogeneous differential equation of 𝒅𝒙
= 𝒈(𝒙,𝒚)
𝑑𝑦 𝑓(𝑥,𝑦)
= is a nonhomogeneous differential equation if 𝑓(𝑥, 𝑦) and 𝑔(𝑥, 𝑦) are nonhomogeneous or
𝑑𝑥 𝑔(𝑥,𝑦)
they have homogeneity with different degrees.
For example:
𝑑𝑦 2𝑥−4𝑦+5
𝑑𝑥
= 𝑥−2𝑦+3
It is a 1st order nonlinear nonhomogeneous differential equation where the functions at numerator and
denominator are nonhomogeneous.
84
To avoid confusion, we use the term “homogeneous/nonhomogeneous” for all the linear ODE case,
𝑑𝑦 𝑓(𝑥,𝑦)
while we use the term “homogeneous/ nonhomogeneous of 𝑑𝑥 = 𝑔(𝑥,𝑦) form” for the 1st order nonlinear
ODE case.
Note: The homogeneity of 2nd and higher order nonlinear ODE is out of scope in this study.
The difference between the solution of algebraic equation and differential equation is shown in table below:
Or, perhaps,
Or, perhaps,
Or, perhaps,
85
The solution of differential equation can be further divided into two types: (a) General Solution (b)
Particular Solution.
𝑑𝑥
For example: The general solution of the differential equation 𝑑𝑡
= −4𝑥 is 𝑥(𝑡) = 𝐴𝑒 −4𝑡 where any
arbitrary constant 𝐴 can satisfy the equation.
𝑑𝑥
For example: Previously, let the general solution of the differential equation = −4𝑥 is 𝑥(𝑡) = 𝐴𝑒 −4𝑡
𝑑𝑡
where any arbitrary constant 𝐴 can satisfy the equation. Given initial condition where 𝑥(0) = 2.5, the
𝑑𝑥
particular solution of the differential equation = −4𝑥 which has the value 2.5 when 𝑡 = 0 is 𝑥(𝑡) =
𝑑𝑡
2.5𝑒 −4𝑡 .
Here, only a specific constant 𝐴 = 2.5 can satisfy the equation.
Note: In fact, general solution & particular solution indicates that there are an infinite number of solutions
to the differential equation unless we are given the specific condition of the problem. The actual solution to
a differential solution is the specific solution that not only satisfies the differential equation, but also satisfies
the given initial conditions.
The particular solution can be obtained from either “Boundary-value problem” or “Initial-value
problem”. It depends on the given specific condition about the value of the solution at a particular point, in
addition to the differential equation.
86
Case (2): Initial-value problem
All conditions are specified at the same value of the independent variable.
𝑑2 𝑥
For example: The particular solution of the differential equation 𝑑𝑡 2
+ 25𝑥 = 0
which has the initial condition 𝑥(0) = 4 & 𝑥 ′ (0) = 8
is 𝑥(𝑡) = 1.6 𝑠𝑖𝑛(5𝑡) + 4𝑐𝑜𝑠(5𝑡) .
Note: For 1st order differential equation, the condition can be treated as initial/ boundary condition. For
higher order differential equation, the distinction becomes obvious.
A solution to an ODE is a function which is differentiable and which satisfies the given equation. This is true
for both explicit and implicit solutions.
The explicit solution is valid in the interval 𝐼: 𝛼 < 𝑡 < 𝛽 I if the following conditions are meet:
𝑑𝑄(𝑡) 𝑑 𝑛−1 𝑄(𝑡)
(i) 𝑄(𝑡), 𝑑𝑡 , … , 𝑑𝑡 𝑛−1 is differentiable
(ii) 𝑄(𝑡) can satisfy the differential equation
87
For example:
𝑑𝑦
Find implicit and explicit solution to the first order differential equation 𝑦 𝑑𝑡 = 𝑡 , 𝑦(2) = −1.
Solution:
𝑑𝑦
𝑦 𝑑𝑡 = 𝑡
𝑦2 𝑡2
>> 2
= 2
+𝐶 [Comment: This is general implicit solution]
(−1)2 (2)2
>> = +𝐶
2 2
3
>> 𝐶 = −
2
Generally, we arrange the implicit solution in the form G(t,y)=0, i.e. 𝑦 2 − 𝑡 2 = −3.
Rearrange the implicit solution and let LHS to be the dependent variable, we get
>> 𝑦1 = +√𝑡 2 − 3
>>𝑦2 = −√𝑡 2 − 3
To check which one is the true solution, reapply the initial condition, 𝑦(2) = −1.
>>> 𝑦1 = +√𝑡 2 − 3 = +√22 − 3 = 1 [Comment: This is not the true explicit solution]
88
Verification of solution:
𝑑𝑦
(i) To verify the 𝑦 2 = 𝑡 2 − 3 is the implicit solution for the differential equation 𝑑𝑡
= 𝑡 . We try to
deduce the differential equation from it.
For example:
𝑦2 = 𝑡2 − 3
𝑑 𝑑
>> 𝑑𝑡 [𝑦 2 ] = 𝑑𝑡 [𝑡 2 − 3]
𝑑𝑦
>> 2𝑦 𝑑𝑡 = 2𝑡
𝑑𝑦
>> 𝑦 =𝑡
𝑑𝑡
𝑑𝑦
∴ Thus, it is proven that 𝑦 2 = 𝑡 2 − 3 is the implicit solution for the 𝑦 𝑑𝑡 = 𝑡.
𝑑𝑦
(ii) To verify the 𝑦2 = −√𝑡 2 − 3 is the explicit solution for the differential equation 𝑑𝑡
= 𝑡 . We
differentiate and substitute it to the equation.
For example:
𝑦2 = −√𝑡 2 − 3
𝑑 𝑑
>> 𝑑𝑡 [𝑦2 ] = 𝑑𝑡 [−√𝑡 2 − 3]
1
𝑑 1
>> 𝑑𝑡 [𝑦2 ] = − 2 (𝑡 2 − 3)2−1 . 2𝑡]
𝑑 𝑡
>> 𝑑𝑡 [𝑦2 ] = −
√(𝑡 2 −3)
𝑑𝑦
Substitute the derivative to the differential equation 𝑦 =𝑡
𝑑𝑡
LHS we get,
𝑑𝑦 𝑡
𝑦 𝑑𝑡 = 𝑦2 (− )
√(𝑡 2 −3)
𝑑𝑦 𝑡
𝑦 = −√𝑡 2 − 3 (− )
𝑑𝑡 √(𝑡 2 −3)
𝑑𝑦
>> 𝑦 =𝑡
𝑑𝑡
∴ Since LHS=RHS, thus it is proven that 𝑦2 = −√𝑡 2 − 3 is the explicit solution for the differential
𝑑𝑦
equation 𝑦 𝑑𝑡 = 𝑡.
(iii) Exercise: Try to prove that 𝑦1 = +√𝑡 2 − 3 is the explicit solution for the differential equation.
89
Note: It will not always be possible to find an explicit solution.
For example:
𝑑𝑦
For differential equation 3𝑦 3 𝑒 3𝑥𝑦 − 1 + (2𝑦𝑒 3𝑥𝑦 + 3𝑥𝑦 2 𝑒 3𝑥𝑦 ) 𝑑𝑥 = 0 , initial condition, 𝑦(0) = 1 . You
suspect that 𝑦 2 𝑒 3𝑥𝑦 − 𝑥 = 1 is the implicit solution to the differential equation. Please verify the solution.
There is no way to rearrange the implicit solution for 𝑦 = 𝑦(𝑥) and get an explicit solution. This mostly
happen for nonlinear case.
90
11.5 STRATEGY TO SOLVE 1 S T ORDER DIFFERENTIAL EQUATION
There are many strategies that have been developed to solve differential equation. We will start with the
most fundamental one. Bear in your mind various strategies can be implemented depends on the
types/forms of differential equation.
First of all, we will start with the strategies to solve 1st order linear differential equation, i.e.
𝑎1 (𝑥)𝑦 ′ + 𝑎0 (𝑥)𝑦 = 𝑔(𝑥). These strategies include
(a) exact differential equation
(b) linear differential equation
(c) separable differential equation
Then we will discuss the strategies used for solving the 1st order nonlinear differential equation, i.e.
𝑎1 (𝑥, 𝑦)(𝑦 ′ )𝐶 + 𝑎0 (𝑥, 𝑦)(𝑦)𝐷 = 𝑔(𝑥, 𝑦). These strategies include
(a) separable differential equation
(b) Bernoulli’s equation
(c) homogeneous differential equation
(d) nonhomogeneous differential equation.
Note: In many cases there are at least one of these strategies which can be used to solve the problem.
Some differential equations are of a form that can be solved readily, it would be useful to be able to
recognize them. For example: we know the derivative for:
91
Identification of exact differential equation would be difficult in most of the time as it involves a lot
of memorization and the identification process become hard, especially when the differential equation is
varies with the original format that you memorize.
For example: Is 𝑥𝑦 ′ + 𝑦 + 4 = 0 an exact differential equation?
Thus, it is crucial for us to know the equation whether it is exact or not before we proceed with some lengthy
process in finding its solution. In general, exact differential equation can be identified by using the following
approach.
𝜕𝑀(𝑥,𝑦) 𝜕𝑁(𝑥,𝑦)
If 𝜕𝑦
= 𝜕𝑥
,
𝜕𝑀(𝑥,𝑦) 𝜕𝑁(𝑥,𝑦)
If 𝜕𝑦
≠ 𝜕𝑥
,
𝑥𝑦 ′ + 𝑦 + 4 = 0
>> 𝑥𝑑𝑦 + (𝑦 + 4)𝑑𝑥 = 0 [Comment: Arrange it in the form of (𝑥, 𝑦)𝑑𝑥 + 𝑁(𝑥, 𝑦)𝑑𝑦 = 0 ]
>> where
𝑀(𝑥, 𝑦) = (𝑦 + 4)
𝑁(𝑥, 𝑦) = 𝑥
𝜕𝑀(𝑥,𝑦) 𝜕𝑁(𝑥,𝑦)
>> Check it with 𝜕𝑦
= 𝜕𝑥
condition.
𝜕𝑀(𝑥,𝑦) 𝜕
LHS: 𝜕𝑦
= 𝜕𝑦 (𝑦 + 4) = 1
𝜕𝑁(𝑥,𝑦) 𝜕
RHS: 𝜕𝑥
= 𝜕𝑥 (𝑥) = 1
𝜕𝑀(𝑥,𝑦) 𝜕𝑁(𝑥,𝑦)
∴ Since 𝜕𝑦
= 𝜕𝑥
is true, we can solve it using exact differential equation.
92
Further solve it using exact approach.
𝑥𝑦 ′ + 𝑦 + 4 = 0
𝑑(𝑥(𝑦+4))
>> 𝑑𝑥
=0
𝑑(𝑥(𝑦+4))
>> ∫ 𝑑𝑥
𝑑𝑥 = ∫ 0𝑑𝑥
>> 𝑥(𝑦 + 4) = 𝐶
𝐶
>> 𝑦 = 𝑥 − 4
Previous section shows that if the differential equation can be in the “exact” form, it can be solved directly.
However, it is undesired to obtain the exact differential equation in most of the time.
𝑥𝑦 ′ − 𝑦 = 0
>> 𝑥𝑑𝑦 − 𝑦𝑑𝑥 = 0 [Comment: Arrange it in the form of (𝑥, 𝑦)𝑑𝑥 + 𝑁(𝑥, 𝑦)𝑑𝑦 = 0 ]
>> where
𝑀(𝑥, 𝑦) = −𝑦
𝑁(𝑥, 𝑦) = 𝑥
𝜕𝑀(𝑥,𝑦) 𝜕𝑁(𝑥,𝑦)
>> Check it with 𝜕𝑦
= 𝜕𝑥
condition.
𝜕𝑀(𝑥,𝑦) 𝜕
LHS: 𝜕𝑦
= 𝜕𝑦 (−𝑦) = −1
𝜕𝑁(𝑥,𝑦) 𝜕
RHS: 𝜕𝑥
= 𝜕𝑥 (𝑥) = 1
𝜕𝑀(𝑥,𝑦) 𝜕𝑁(𝑥,𝑦)
∴ Since 𝜕𝑦
= 𝜕𝑥
is not true, we cannot solve it using exact differential equation.
Note: It can’t be solved as it is in “non exact form”. However, there is a method used to convert it to exact
form. It is a strategy involving Integrating Factor and Linear Differential Equation as shown in the
following steps:
93
Procedure to solve 1st linear ODE problem using linear differential equation
𝑑 𝑑𝑦 𝑑
Prove: 𝑑𝑥 (𝑦. 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 ) = (𝑑𝑥 ) 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 + 𝑦 (𝑑𝑥 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 )
𝑑𝑦 𝑑
= ( ) 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 + 𝑦 ( 𝑝(𝑥)𝑑𝑥) 𝑒 ∫ 𝑝(𝑥)𝑑𝑥
𝑑𝑥 𝑑𝑥
𝑑𝑦
= ( ) 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 + 𝑦(𝑝(𝑥))𝑒 ∫ 𝑝(𝑥)𝑑𝑥
𝑑𝑥
𝑑𝑦
= 𝐼𝐹. (𝑑𝑥 ) + 𝐼𝐹. (𝑝(𝑥)𝑦) [proven]
Step 5:
𝑑
Integrate ∫ 𝑦(𝐼𝐹)𝑑𝑥 = ∫ 𝐼𝐹(𝑞(𝑥))
𝑑𝑥
94
Continue with the previous problem:
Solve 𝑥𝑦 ′ − 𝑦 = 0 [where 𝑥 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒 and 𝑦 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒]
𝑑𝑦 1 𝑑𝑦
>> 𝑑𝑥 − 𝑥 𝑦 = 0 [Step 1- linear form of 𝑑𝑥
+ 𝑝(𝑥)𝑦 = 𝑞(𝑥) ]
where
1
𝑝(𝑥) = − 𝑥
𝑞(𝑥) = 0
1
>> The integrating factor, 𝐼𝐹 = 𝑒 ∫ −𝑥𝑑𝑥 = 𝑒 −𝑙𝑛𝑥 = 𝑥 −1 [Step 2- IF=𝑒 ∫ 𝑝(𝑥)𝑑𝑥 ]
𝑑𝑦 1
>> 𝑥 −1 ( ) − 𝑥 −1 ( 𝑦) = 𝑥 −1 (0) [Step 3- Multiply]
𝑑𝑥 𝑥
𝑑𝑦 1
>>𝑥 −1 ( ) − ( 𝑦) =0
𝑑𝑥 𝑥2
𝑑
>> 𝑑𝑥 (𝑦. 𝑥 −1 ) = 0 [Step 4- Exact]
𝑑
where LHS= (𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒. 𝐼𝐹)
𝑑𝑥
𝑑
>>∫ 𝑑𝑥 (𝑥 −1 𝑦)𝑑𝑥 = ∫ 0𝑑𝑥 [Step 5- Integrate]
>> 𝑥 −1 𝑦 = 𝐶
>> 𝑦 = 𝐶𝑥
∴ The solution of the first order linear ODE problem 𝑥𝑦 ′ − 𝑦 = 0 is 𝑦 = 𝐶𝑥, where 𝐶 = arbitrary constant.
95
𝑑𝑥
Example: Solve the differential equation 𝑦
+ 2𝑥𝑑𝑦 = 0 by using Linear Differential Equation.
Solution:
𝑑𝑥
𝑦
+ 2𝑥𝑑𝑦 = 0
1 𝑑𝑦
>> 𝑦
+ 2𝑥 𝑑𝑥 = 0 [Comment: This is a nonlinear ODE; Rearrange it to linear ODE]
where
𝑝(𝑦) = 2𝑦
𝑞(𝑦) = 0
2
>> The integrating factor, 𝐼𝐹 = 𝑒 ∫ 2𝑦𝑑𝑥 = 𝑒 𝑦 [Step 2- IF]
2 𝑑𝑥 2
>> 𝑒 𝑦 𝑑𝑦
+ 𝑒 𝑦 (2𝑦)𝑥 = 0 [Step 3- Multiply]
𝑑 2
>> 𝑑𝑦 (𝑥. 𝑒 𝑦 ) = 0 [Step 4- Exact]
𝑑 2
>> ∫ 𝑑𝑦 (𝑥. 𝑒 𝑦 )𝑑𝑦 = ∫ 0𝑑𝑦
2
>>(𝑥. 𝑒 𝑦 ) = 𝐶
2
∴ 𝑥 = 𝐶𝑒 −𝑦 [Step 5- Integrate]
96
𝑑𝑦 𝑦
Example: Solve the initial value problem 𝑑𝑥 = 2𝑥+3𝑦2 −2 , 𝑦(1) = 1
Solution:
𝑑𝑦 𝑦
𝑑𝑥
= 2𝑥+3𝑦2 −2
𝑑𝑦 𝑦
>> − =0 [Comment: This is a nonlinear ODE; Rearrange it to linear ODE]
𝑑𝑥 2𝑥+3𝑦 2 −2
𝑑𝑥 2 3𝑦 2 −2
>> 𝑑𝑦 − (𝑦) 𝑥 = 𝑦
[Step 1- linear form]
2
∫ −(𝑦)𝑑𝑥
>> The integrating factor, 𝐼𝐹 = 𝑒 = 𝑒 −2𝑙𝑛𝑦 = 𝑦 −2 [Step 2- IF]
𝑑𝑥 2 3𝑦 2 −2 3𝑦 2 −2
>> 𝑦 −2 𝑑𝑦 − 𝑦 −2 (𝑦) 𝑥 = 𝑦 −2 ( 𝑦
) =( 𝑦3
) [Step 3- Multiply]
𝑑 3𝑦 2 −2
>> (𝑥. 𝑦 −2 ) =( ) [Step 4- Exact]
𝑑𝑦 𝑦3
𝑑 3𝑦 2 −2
>> ∫ 𝑑𝑦 (𝑥. 𝑦 −2 )𝑑𝑦 = ∫ ( 𝑦3
) 𝑑𝑦 [Step 5- Integrate]
3 −2
>> 𝑥. 𝑦 −2 = ∫ ( ) 𝑑𝑦 + ∫ ( 3 ) 𝑑𝑦
𝑦 𝑦
∴ 𝑥 = 3𝑦 2 𝑙𝑛𝑦 + 1
97
In this study, it shows that integrating factor 𝐼𝐹 = 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 is very useful in solving 1st order linear
ODE. For your additional knowledge, there are other types of integrating factors to solve various ODE
problems, as shown in the table below:
Note: This is for your extra knowledge and other types of integrating factors are not included in this study.
Then, we can solve the problem by integrating both sides. This is particularly useful in solving certain linear
& nonlinear differential equation which is separable.
98
𝑑𝑦
For example: Solve the equation 9𝑦 𝑑𝑥 − 4𝑥 = 0.
Solution:
Recognize that this is a nonlinear equation that is difficult to be solved by using Linear or Exact Differential
Equations. However, it can be solved easily by using the Separable Differential Equation.
𝑑𝑦
9𝑦 𝑑𝑥 − 4𝑥 = 0 [Comment- Nonlinear form]
>> 𝑡𝑎𝑛−1 𝑦 = 𝑥 + 𝐶
∴ 𝑦 = 𝑡𝑎𝑛(𝑥 + 𝐶)
𝑑𝑦 𝑦+1
For example: Solve the equation 𝑑𝑥 = 𝑥−4 , 𝑦(6) = 0
Solution:
𝑑𝑦 𝑦+1
𝑑𝑥
= 𝑥−4 [Comment- Nonlinear form]
1 1
>>∫ 𝑦+1 𝑑𝑦 = ∫ 𝑥−4 𝑑𝑥
There are some cases cannot be solved by the previous strategies due to its non-linear and non-separable
characteristics, however, if it is in the form of Bernoulli’s equation, it can be converted to the linear
differential equation and hence solved by using the Integrating Factor strategy.
𝑑𝑦
For 𝑛 = 0, 𝑑𝑥
+ 𝑝(𝑥)𝑦 = 𝑟(𝑥) [linear differential equation]
𝑑𝑦
For 𝑛 = 1, 𝑑𝑥
+ (𝑝(𝑥) − 𝑟(𝑥))𝑦 = 0 [linear differential equation]
𝑑𝑦
For 𝑛 > 1, 𝑑𝑥
+ 𝑝(𝑥)𝑦 = 𝑟(𝑥)𝑦 𝑛 [nonlinear differential equation]
To convert the nonlinear Bernoulli’s equation into linear form, we need two important properties:
[i] Let 𝑣(𝑥) = {𝑦(𝑥)}1−𝑛
where 𝑦 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒 & 𝑥 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
1 𝑑𝑣(𝑥) 𝑑𝑦(𝑥)
[ii] The derivative, (1−𝑛) 𝑑𝑥
= {𝑦(𝑥)}−𝑛 𝑑𝑥
Prove:
𝑑𝑣(𝑥) 𝑑𝑦(𝑥)
= (1 − 𝑛){𝑦(𝑥)}(1−𝑛)−1
𝑑𝑥 𝑑𝑥
1 𝑑𝑣(𝑥) 𝑑𝑦(𝑥)
>> (1−𝑛) 𝑑𝑥
= {𝑦(𝑥)}−𝑛 𝑑𝑥
𝑑𝑦
Divide by 𝑦 𝑛 and we get 𝑦 −𝑛 𝑑𝑥 + 𝑝(𝑥)𝑦1−𝑛 = 𝑟(𝑥) [Step 1: Divide by 𝑦 𝑛 ]
Hence, the 1st order linear differential equation can be solved using the previous strategy. Once you obtain
the solution 𝑣. Back substitute it into eqn. 𝑣(𝑥) = {𝑦(𝑥)}1−𝑛 to obtain the solution 𝑦.
𝑑𝑦 3
For example: Solve the equation 𝑑𝑥 + 3𝑥 2 𝑦 = 𝑥𝑒 𝑥 𝑦 2
Solution:
𝑑𝑦 3
+ 3𝑥 2 𝑦 = 𝑥𝑒 𝑥 𝑦 2
𝑑𝑥
[Step 1: Divide by 𝑦 𝑛 ]
𝑑𝑦 3
>> 𝑦 −2 𝑑𝑥 + 3𝑥 2 𝑦 −1 = 𝑥𝑒 𝑥
101
𝑑𝑦 3
𝑦 −2 + 3𝑥 2 𝑦 −1 = 𝑥𝑒 𝑥
𝑑𝑥
𝑑𝑣(𝑥) 3
>> − 𝑑𝑥
+ 3𝑥 2 𝑣(𝑥) = 𝑥𝑒 𝑥 [Comment: linear Bernoulli’s equation]
where
𝑣 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑥 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑝(𝑥)=−3𝑥 2
3
𝑟(𝑥) = −𝑥𝑒 𝑥
2 𝑑𝑥 3
>> The integrating factor, 𝐼𝐹 = 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 = 𝑒 ∫ −3𝑥 = 𝑒 −𝑥 [Step 2- IF]
3 𝑑𝑣(𝑥) 3 3 3
>> 𝑒 −𝑥 𝑑𝑥
− 𝑒 −𝑥 (3𝑥 2 𝑣(𝑥)) = 𝑒 −𝑥 (−𝑥𝑒 𝑥 ) [Step 3- Multiply]
3 𝑑𝑣(𝑥) 3
>> 𝑒 −𝑥 𝑑𝑥
− 𝑒 −𝑥 (3𝑥 2 𝑣(𝑥)) = −𝑥
𝑑 −𝑥 3
>> 𝑑𝑥
(𝑒 . 𝑣(𝑥)) = −𝑥 [Step 4- Exact]
3
where 𝐼𝐹 = 𝑒 −𝑥 & 𝑣 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑑 3
>> ∫ 𝑑𝑥 (𝑒 −𝑥 . 𝑣(𝑥)) 𝑑𝑥 = ∫(−𝑥)𝑑𝑥 [Step 5- Integrate]
3 −𝑥 2
>> 𝑒 −𝑥 . 𝑣(𝑥) = 2
+𝐶
102
𝑑𝑦
For example: Solve the equation 3 𝑑𝑥 + 𝑦 = (1 − 2𝑥)𝑦 4 .
Solution:
𝑑𝑦
3 + 𝑦 = (1 − 2𝑥)𝑦 4
𝑑𝑥
𝑑𝑦
Recognize that it is nonlinear Bernoulli eqn. 𝑑𝑥 + 𝑝(𝑥)𝑦 = 𝑟(𝑥)𝑦 𝑛 ,
1 (1−2𝑥)
With 𝑝(𝑥) = , 𝑟(𝑥) = ,𝑛 =4
3 3
[Step 1: Divide by 𝑦 𝑛 ]
𝑑𝑦 𝑦 (1−2𝑥) 4
𝑑𝑥
+3= 3
𝑦
𝑑𝑦 𝑦 −3 (1−2𝑥)
>> 𝑦 −4 𝑑𝑥 + 3
= 3
𝑑𝑦 𝑦 −3 (1−2𝑥)
𝑦 −4 𝑑𝑥 + 3
= 3
1 𝑑𝑣(𝑥) 𝑣(𝑥) (1−2𝑥)
>> −3 𝑑𝑥
+ 3
= 3
𝑑𝑣(𝑥)
>> − 𝑑𝑥
+ 𝑣(𝑥) = (1 − 2𝑥) [Comment: linear Bernoulli’s equation]
where
𝑣 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑥 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑝(𝑥)=−1
𝑟(𝑥) = (2𝑥 − 1)
103
>> The integrating factor, 𝐼𝐹 = 𝑒 ∫ 𝑝(𝑥)𝑑𝑥 = 𝑒 ∫ −1𝑑𝑥 = 𝑒 −𝑥 [Step b- IF]
𝑑𝑣(𝑥)
>> 𝑒 −𝑥 𝑑𝑥
− 𝑒 −𝑥 𝑣(𝑥) = 𝑒 −𝑥 (2𝑥 − 1) [Step c- Multiply]
𝑑
>> 𝑑𝑥
(𝑒 −𝑥 . 𝑣(𝑥)) = 𝑒 −𝑥 (2𝑥 − 1) [Step d- Exact]
where 𝐼𝐹 = 𝑒 −𝑥 & 𝑣 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑑
>> ∫
𝑑𝑥
(𝑒 −𝑥 . 𝑣(𝑥))𝑑𝑥 = ∫ 𝑒 −𝑥 (2𝑥 − 1)𝑑𝑥 [Step e- Integrate]
3 1
∴ 𝑦(𝑥) = √−2𝑥−1+𝐶𝑒 𝑥 , where 𝐶 = arbitrary constant.
104
11.5.5 DIFFERENTIAL EQUATION OF HOMOGENEOUS 𝑑𝑦/𝑑𝑥 = 𝑓(𝑥, 𝑦)/𝑔(𝑥, 𝑦) FORM
𝑑𝑦 𝑓(𝑥,𝑦)
Previously, we use the 𝑑𝑥 = 𝑔(𝑥,𝑦) form to check the homogeneity of 1st order nonlinear ODE. This section
𝑑𝑦 𝑓(𝑥,𝑦)
introduces the strategy used to solve the homogeneous 𝑑𝑥 = 𝑔(𝑥,𝑦).
𝑑𝑦
For example: Solve 2𝑥𝑦 − 𝑦 2 = −𝑥 2
𝑑𝑥
𝑑𝑦 𝑓(𝑥,𝑦)
Rearrange it to the form of = , we get
𝑑𝑥 𝑔(𝑥,𝑦)
𝑑𝑦 𝑦 2 −𝑥 2
>> =
𝑑𝑥 2𝑥𝑦
𝑑𝑦 𝑓(𝑥,𝑦)
The 1st order nonlinear ODE with homogeneous 𝑑𝑥 = 𝑔(𝑥,𝑦) form can be solved by using substitution
method as shown in the table below. In general, this form is very useful to convert the non-separable
differential equation into separable differential equation.
105
𝑑𝑦 𝑓(𝑥,𝑦)
To convert the 1st order nonlinear ODE with homogeneous 𝑑𝑥 = 𝑔(𝑥,𝑦) form into separable form, we need
two important properties:
𝑦
[i] Let 𝑣(𝑥) = 𝑥 or 𝑦 = 𝑣𝑥
where 𝑦 = 𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒 & 𝑥 = 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 𝑣𝑎𝑟𝑖𝑎𝑏𝑙𝑒
𝑑𝑦 𝑑
[ii] The derivative, 𝑑𝑥 = 𝑥 𝑑𝑥 [𝑣] + 𝑣
𝑑𝑦
For example: Solve 2𝑥𝑦 − 𝑦 2 = −𝑥 2
𝑑𝑥
Solution:
𝑑𝑦
2𝑥𝑦 𝑑𝑥 − 𝑦 2 = −𝑥 2
𝑑𝑦 𝑦 2 −𝑥 2
>> 𝑑𝑥 = 2𝑥𝑦
𝑑𝑦 𝑓(𝑥,𝑦)
[Comment: 1st order nonlinear ODE with homogeneous 𝑑𝑥 = 𝑔(𝑥,𝑦) form; non-separable form]
𝑦 = 𝑣𝑥
𝑑𝑦 𝑑
=𝑥 [𝑣] + 𝑣
𝑑𝑥 𝑑𝑥
𝑑𝑦 (𝑣𝑥)2 −𝑥 2 𝑑
>> = =𝑥 [𝑣] + 𝑣
𝑑𝑥 2𝑥(𝑣𝑥) 𝑑𝑥
(𝑣)2 −1 𝑑𝑣
>> −𝑣 = 𝑥 [Comment: Solve using separable differential equation]
2(𝑣) 𝑑𝑥
1 1
>> 𝑥 𝑑𝑥 = (𝑣)2−1 𝑑𝑣
−𝑣
2(𝑣)
1 1
>> 𝑥 𝑑𝑥 = 𝑣2 −1−2𝑣2
𝑑𝑣
2(𝑣)
1 2𝑣
>> 𝑥 𝑑𝑥 = −𝑣 2 −1 𝑑𝑣 [Step a- Separable form]
1 2𝑣
>>∫ 𝑥 𝑑𝑥 = ∫ −𝑣2 −1 𝑑𝑣 [Step b- Integrate both sides]
106
Step 2: Back Substitution
𝑦 2
>> 𝑙𝑛|𝑥| = −𝑙𝑛 |− (𝑥 ) − 1| + 𝐶
𝑦 2
>> 𝑙𝑛|𝑥| + 𝑙𝑛 |− (𝑥 ) − 1| = +𝐶
𝑦2
>> 𝑙𝑛 |− − 𝑥| = +𝐶
𝑥
𝑦2
>> − 𝑥
− 𝑥 = 𝑒𝐶
>> 𝑦 2 = −𝑥 2 − 𝑥𝑒 𝐶
∴ 𝑦 = ±√−𝑥 2 − 𝑥𝑒 𝐶
Case I Case II
𝑎1 𝑏1 𝑎 𝑏1
| |=0 | 1 |≠0
𝑎2 𝑏2 𝑎2 𝑏2
Let 𝑣 = 𝑎1 𝑥 + 𝑏1 𝑦 𝑥 =𝑋+ℎ
Let
𝑦 = 𝑌+𝑘
where the pair constant (ℎ, 𝑘) can be obtained by
solving the simultaneous equations:
𝑎1 ℎ + 𝑏1 𝑘 + 𝑐1 = 0
𝑎2 ℎ + 𝑏2 𝑘 + 𝑐2 = 0
𝑑𝑦 𝑓(𝑥,𝑦)
The substitution converts the nonhomogeneous = to homogeneous form and hence it can be
𝑑𝑥 𝑔(𝑥,𝑦)
solved by method introduced in section 11.5.5.
𝑑𝑦 𝑓(𝑥,𝑦)
The nonlinear nonhomogeneous 𝑑𝑥 = 𝑔(𝑥,𝑦) has a lot of branches and it may require various types
of substitution to convert it to homogeneous form. The topic is wide and thus it is not covered in this study.
107
SOLUTION TO HOMOGENEOUS & NON-
HOMOGENEOUS
WEEK 12: SOLUTION TO HOMOGENEOUS & NON-HOMOGENEOUS
12.1 INTRODUCTION TO 2 ND ORDER DIFFERENTIAL EQUATION
As discussed earlier, the differential equation is important, especially in mathematical modelling for
engineering application. Previously, we have discussed several strategies and methods to solve the 1st order
differential equation. However, the 1st order differential equation is only able to model certain engineering
problem.
For example,
However, this differential equation is insufficient to describe a 𝑅𝐿𝐶 circuit as illustrated below:
where the RLC electric circuit consists a generator (𝐸 volt), a resistance (𝑅 ohms), a inductance (𝐿 Henry), a
capacitor (C capacitance).
We know that the voltage loss through a resistor, capacitor and inductor is proportional to the current,
charge and rate of charge change respectively, where
Capacitor 1
𝛥𝑉(𝑡) = (𝑉𝐹 − 𝑉𝐵 ) = 𝐶 𝑞(𝑡)
𝑑𝑞(𝑡) 1 𝑑 2 𝑞(𝑡)
>> 𝑅 𝑑𝑡
+ 𝐶 𝑞(𝑡) + 𝐿 𝑑𝑡 2
− 𝐸(𝑡) = 0.
𝑑 2 𝑞(𝑡) 𝑑𝑞(𝑡) 1
>> 𝐿 +𝑅 + 𝑞(𝑡) = 𝐸(𝑡).
𝑑𝑡 2 𝑑𝑡 𝐶
In this chapter, we will discuss about the strategies used to solve the 2nd order differential equation.
Previously, we know that 2nd order linear differential equation can be categorized into two forms:
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
(i) Homogeneous (i.e. 𝑎 𝑑𝑥 2
+𝑏 𝑑𝑥
+ 𝑐𝑦(𝑥) = 0.)
Solution: Known as complementary solution, 𝑦𝑐
Example:
𝑑2 𝑦 𝑑𝑦(𝑥)
The solution to 𝑑𝑥 2 − 4 𝑑𝑥
+ 3𝑦(𝑥) = 0 is 𝑦𝑐 = 𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
(ii) Nonhomogeneous (i.e. 𝑎 𝑑𝑥 2
+𝑏 𝑑𝑥
+ 𝑐𝑦(𝑥) = 𝑟(𝑥)).
Solution: Combination of complementary solution, 𝑦𝑐 and particular solution, 𝑦𝑝
Example:
𝑑2 𝑦 𝑑𝑦(𝑥) 2
The solution to 𝑑𝑥 2 − 4 𝑑𝑥
+ 3𝑦(𝑥) = 10𝑒 −2𝑥 is 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥 + 3 𝑒 −2𝑥
109
Note: We assume 𝑦 = dependent variable and 𝑥 = independent variable in the above example. It is
worthwhile to mention that the methods to obtain the yc is same for both homogeneous and
nonhomogeneous cases. Just that we have an additional solution y p for nonhomogeneous case.
Before we proceed to the strategy in solving the 2nd order differential equation, we should know about the
theory of linearity principle and linearly dependency.
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
It is given that the solution to 𝑑𝑥 2
−4 𝑑𝑥
+ 3𝑦(𝑥) = 0 is 𝑦𝑐 = 𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥
Substitute to LHS
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
𝑑𝑥 2
−4 𝑑𝑥
+ 3𝑦(𝑥) = 9𝑒 3𝑥 − 4(3𝑒 3𝑥 ) + 3𝑒 3𝑥 = 0
>> LHS = RHS=0
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
∴ 𝑦1 = 𝑒 3𝑥 is proven to be the solution of 𝑑𝑥 2
−4 𝑑𝑥
+ 3𝑦(𝑥) = 0
110
Substitute to LHS
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
−4 + 3𝑦(𝑥) = 𝑒 𝑥 − 4(𝑒 𝑥 ) + 3𝑒 𝑥 = 0
𝑑𝑥 2 𝑑𝑥
Substitute to LHS
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
𝑑𝑥 2
−4 𝑑𝑥
+ 3𝑦(𝑥) = (9𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥 ) − 4(3𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥 ) + 3(𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥 ) = 0
>> LHS = RHS=0
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
∴ 𝑦𝑐 = 𝑐1 𝑒 3𝑥 + 𝑐2 𝑒 𝑥 is proven to be the solution of 𝑑𝑥 2
−4 𝑑𝑥
+ 3𝑦(𝑥) = 0
From the verification, it was shown that all the three solutions satisfy the ODE equation and thus they are
true solutions. The complete solution is formed according to Linearity Principle/ Principle of Superposition
as follows:
𝑑 2 𝑦(𝑥) 𝑑𝑦(𝑥)
So, the solution of the 2nd order differential equation ( 𝑑𝑥 2
−4 𝑑𝑥
+ 3𝑦(𝑥) = 0) is not purely
3𝑥 𝑥
𝑦1 = 𝑒 or 𝑦2 = 𝑒 . If a 2 order linear differential equation is encountered, the complete solution will be
nd
111
In conclusion, the general complementary solution of the homogeneous linear differential
equation is equal to
𝑦 = 𝑐1 𝑦1 (𝑥) + 𝑐2 𝑦2 (𝑥)
where 𝑦1 (𝑥) & 𝑦2 (𝑥) are also the solution for the equation and they are linearly independent to each other.
Frankly speaking, if two solutions are linearly dependent, it means these two solutions are
redundant and hence they do not represent two independent solutions instead of one. Thus, linearly
independent solutions of 𝑦1 (𝑥) & 𝑦2 (𝑥) are desired. We can use the following two methods to check
whether two solutions are linearly independent to each other or not: (a) Linear Dependency Theorem (b)
Wronskian Method.
Or in other words,
𝑦1 (𝑥) & 𝑦2 (𝑥) are proportional to each other
Or in other words,
𝑦1 (𝑥) & 𝑦2 (𝑥) are not proportional to each other
112
For example: If 𝑦1 = 𝑒 𝑥 & 𝑦2 = 𝑒 3𝑥 are linearly independent, 𝑐1 𝑒 𝑥 + 𝑐2 (2𝑒 𝑥 ) = 0 , only
when 𝑐1 & 𝑐2 = 0.
Check if 𝑦1 = 𝑒 𝑥 & 𝑦2 =
𝑒 3𝑥 are linearly
dependent or not.
It was found that only when 𝑐1 = 0 & 𝑐2 = 0,
𝑐1 𝑒 𝑥 + 𝑐2 (𝑒 3𝑥 ) = 0𝑒 𝑥 + 0𝑒 3𝑥 = 0
Note:
If linearly dependent solutions are obtained,
i.e. 𝑦1 & 𝑦2 are linearly dependent,
we do not obtained a complete solution of
𝑦 = 𝑐1 𝑦1 (𝑥) + 𝑐2 𝑦2 (𝑥).
Method 2
(Wronskian, 𝑊(𝑦1 , 𝑦2 ))
(i) Linearly dependent 𝑦1 (𝑥) & 𝑦2 (𝑥) are linearly dependent if
[Undesired solutions for ODE] 𝑦1 𝑦2
𝑑𝑦
𝑊(𝑦1 , 𝑦2 ) = | 1 𝑑𝑦 2| = 0
𝑑𝑥 𝑑𝑥
For example: 𝑦1 = 𝑒 𝑥
Check if 𝑦1 = 𝑒 𝑥 & 𝑦2 = 2𝑒 𝑥 >> 𝑑𝑦1 = 𝑒 𝑥
𝑑𝑥
are linearly dependent or not.
𝑦2 = 2𝑒 𝑥
𝑑𝑦
>> 𝑑𝑥2 = 2𝑒 𝑥
𝑦1 𝑦2
𝑒𝑥 2𝑒 𝑥
𝑊(𝑦1 , 𝑦2 ) = |𝑑𝑦1 𝑑𝑦2 | =| 𝑥 | = 𝑒 𝑥 (2𝑒 𝑥 ) − 𝑒 𝑥 (2𝑒 𝑥 ) = 0
𝑑𝑥 𝑑𝑥 𝑒 2𝑒 𝑥
113
(ii) Linearly independent 𝑦1 (𝑥) & 𝑦2 (𝑥) are linearly independent if
[Desired solutions of ODE] 𝑦1 𝑦2
𝑊(𝑦1 , 𝑦2 ) = |𝑑𝑦1 𝑑𝑦2 | ≠ 0
𝑑𝑥 𝑑𝑥
For example: 𝑦1 = 𝑒 𝑥
𝑑𝑦1
Check if 𝑦1 = 𝑒 𝑥 & 𝑦2 = 𝑒 3𝑥 >> = 𝑒𝑥
𝑑𝑥
are linearly dependent or
not.
𝑦2 = 𝑒 3𝑥
𝑑𝑦2
>> = 3𝑒 3𝑥
𝑑𝑥
𝑦1 𝑦2 𝑥
𝑊(𝑦1 , 𝑦2 ) = |𝑑𝑦1 𝑑𝑦2 | = |𝑒 𝑥 𝑒 3𝑥 | = 𝑒 𝑥 (3𝑒 3𝑥 ) − 𝑒 𝑥 (𝑒 3𝑥 ) = 2𝑒 4𝑥
𝑑𝑥 𝑑𝑥 𝑒 3𝑒 3𝑥
Since 𝑊(𝑦1 , 𝑦2 ) ≠ 0
∴ Thus, 𝑦1 = 𝑒 𝑥 & 𝑦2 = 𝑒 3𝑥 are linearly independent.
Note:
If linearly dependent solutions are obtained,
i.e. 𝑦1 & 𝑦2 are linearly dependent,
we do not obtained a complete solution of
𝑦 = 𝑐1 𝑦1 (𝑥) + 𝑐2 𝑦2 (𝑥).
114
12.4 STRATEGY TO SOLVE 2 N D ORDER DIFFERENTIAL EQUATION
In this study, we will discuss several strategies to solve the 2nd order linear differential equations that are
given in the following form:
(iii) Nonhomogeneous linear differential equation with constant coefficients 𝑎, 𝑏, 𝑐, in the form of
𝑟(𝑥) = 𝑒 𝛼𝑥 𝑃𝑛 (𝑥)
𝑑2 𝑦 𝑑𝑦
𝑎 +𝑏 + 𝑐𝑦 = 𝑟(𝑥) [Strategy: Method of underdetermined coefficients]
𝑑𝑥 2 𝑑𝑥
Over the years, scientist and engineer have found that the solution to 2nd order differential equation:
𝑑2 𝑦 𝑑𝑦
Solution of 𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 0
to be 𝑦(𝑥) = 𝑒 𝑚𝑥 ≠ 0
However, this is not a complete solution because 2nd order problem should have 2 linearly independent
solutions. To find the complete solution to the problem, we therefore follow the strategy utilizing the
characteristic/auxiliary equation.
𝑑2 𝑦 𝑑𝑦
To solve 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 0 ,
115
Prove: To obtain the characteristic equation: 𝑎𝑚2 + 𝑏𝑚 + 𝑐 = 0
𝑑2 𝑦 𝑑𝑦
𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 0
Assume solution: 𝑦(𝑥) = 𝑒 𝑚𝑥 to be solution of 2nd order ODE
>> Its derivative: 𝑦 ′ = 𝑚𝑒 𝑚𝑥 ; 𝑦′′ = 𝑚2 𝑒 𝑚𝑥
𝑑2 𝑦 𝑑𝑦
Substitute into equation 𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 0 , we get
Recall: Complex conjugate has same magnitude but opposite sign for the imaginary part. For example, the
complex conjugate for a complex number 𝑚1 = (5 + 6𝑖) is 𝑚2 = (5 − 6𝑖).
The summary of the complete solution are listed below. The detail description will be provided next.
Type of (a) Real and distinct (b) A pair of complex conjugates (c) Repeated real root
Roots roots roots 𝑚 = 𝑚1 = 𝑚2
𝑚1 & 𝑚2 𝑚1 = 𝑚 + 𝑖𝛽 & 𝑚2 = 𝑚 − 𝑖𝛽
Indicator 𝑏 2 − 4𝑎𝑐 > 0 𝑏 2 − 4𝑎𝑐 < 0 𝑏 2 − 4𝑎𝑐 = 0
Complete
𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑥𝑒 𝑚2 𝑥 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑥𝑒 𝑚2 𝑥
solution
Or
116
Case (a): Real and distinct roots
𝑚1 ≠ 𝑚2
Characteristic equation:
𝑎𝑚2 + 𝑏𝑚 + 𝑐 = 0
Indicator:
>> 𝑏 2 − 4𝑎𝑐 > 0
Comment: If 𝑏 2 − 4𝑎𝑐 is greater than 0, it indicates that the roots, 𝑚1 & 𝑚2 are real and distinct.
Complete solution:
>> 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥
where 𝑚1 ≠ 𝑚2
𝑑2 𝑦 𝑑𝑦
𝑑𝑥 2
+ 2 𝑑𝑥 − 3𝑦 = 0
>> 𝑚1 = 1, 𝑚3 = −3
Complete solution:
∴ 𝑦(𝑥) = 𝑐1 𝑒 𝑥 + 𝑐2 𝑒 −3𝑥
117
Case (b): A pair of complex conjugates roots
𝑚1 = 𝑚 + 𝑖𝛽 & 𝑚2 = 𝑚 − 𝑖𝛽
Characteristic equation:
𝑎𝑚2 + 𝑏𝑚 + 𝑐 = 0
Indicator:
>> 𝑏 2 − 4𝑎𝑐 < 0
Comment: If 𝑏 2 − 4𝑎𝑐 is less than 0, it indicates that the roots, 𝑚1 & 𝑚2 are a pair of complex conjugates.
Complete solution:
>> 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥
where 𝑚1 ≠ 𝑚2 ;
𝑚1 = 𝑚 + 𝑖𝛽 & 𝑚2 = 𝑚 − 𝑖𝛽;
𝑖 = √−1 = 𝑖𝑚𝑎𝑔𝑖𝑛𝑎𝑟𝑦
or
>> 𝑦(𝑥) = 𝑒 𝑚𝑥 (𝐴𝑐𝑜𝑠𝛽𝑥 + 𝐵𝑠𝑖𝑛𝛽𝑥)
where 𝐴 = 𝑐1 + 𝑐2 ;
𝐵 = 𝑖(𝑐1 − 𝑐2 )
Note: In this case, the complete solution can be either 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥 and 𝑦(𝑥) =
𝑒 𝑚𝑥 (𝐴𝑐𝑜𝑠𝛽𝑥 + 𝐵𝑠𝑖𝑛𝛽𝑥). Both are the same equation but in different format. See the proof below.
𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥
>> 𝑦(𝑥) = 𝑐1 𝑒 (𝑚+𝑖𝛽)𝑥 + 𝑐2 𝑒 (𝑚−𝑖𝛽)𝑥
= 𝑐1 𝑒 (𝑚)𝑥 𝑒 (𝑖𝛽)𝑥 + 𝑐2 𝑒 (𝑚)𝑥 𝑒 (−𝑖𝛽)𝑥
= 𝑒 (𝑚)𝑥 ( 𝑐1 𝑒 (𝑖𝛽)𝑥 + 𝑐2 𝑒 (−𝑖𝛽)𝑥 )
118
Example of the case of a pair of complex conjugates roots:
𝑑2 𝑦 𝑑𝑦
4 + 16 + 17𝑦 = 0
𝑑𝑥 2 𝑑𝑥
𝑚2 + 4𝑚 + 17⁄4 = 0
Complete solution:
1 1
∴ 𝑦(𝑥) = 𝑐1 𝑒 (−2+2𝑖)𝑥 + 𝑐2 𝑒 (−2−2𝑖)𝑥
Or
1 1
𝑦(𝑥) = 𝑒 −2𝑥 (𝐴𝑐𝑜𝑠 𝑥 + 𝐵𝑠𝑖𝑛 𝑥)
2 2
119
Case (c): Repeated real root
𝑚 = 𝑚1 = 𝑚2
Characteristic equation:
𝑎𝑚2 + 𝑏𝑚 + 𝑐 = 0
Indicator:
>> 𝑏 2 − 4𝑎𝑐 = 0
Comment: If 𝑏 2 − 4𝑎𝑐 is equal to 0, it indicates that the roots, 𝑚1 & 𝑚2 are repeated real root.
Complete solution:
Previous solution is not valid here
>> 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥
where 𝑚1 = 𝑚2 and
thus 𝑐1 𝑒 𝑚1 𝑥 & 𝑐2 𝑒 𝑚2 𝑥 are linearly dependent solution (undesired situation) in this case.
𝑑2 𝑦 𝑑𝑦
𝑑𝑥 2
− 4 𝑑𝑥 + 4𝑦 = 0
𝑚2 − 4𝑚 + 4 = 0
>> (𝑚 − 2)(𝑚 − 2) = 0
>> 𝑚1 = 2, 𝑚2 = 2
Complete solution:
∴ 𝑦(𝑥) = 𝑐1 𝑒 2𝑥 + 𝑐2 𝑥𝑒 2𝑥
120
Overall comment:
1. Suppose that the roots of the characteristic equation are 𝑚1 & 𝑚2 , then 𝑒 𝑚1 𝑥 & 𝑒 𝑚2 𝑥 are the
solution of the differential equation.
𝑑2 𝑦 𝑑𝑦
2. Since 𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 0 is a linear homogeneous equation, by using the linear independency
and linearity principle, the general solution must be
(i) 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥 for real and distinct root,
𝑚1 ≠ 𝑚2 .
(ii) 𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑒 𝑚2 𝑥 for a pair of complex conjugate roots, 𝑚1 = 𝑚 + 𝑖𝛽 & 𝑚2 = 𝑚 −
𝑖𝛽.
(iii) However, if there is repeated real roots, 𝑚 = 𝑚1 = 𝑚2 , we get linearly dependent solution
𝑦(𝑥) = 𝑐1 𝑒 𝑚𝑥 + 𝑐2 𝑒 𝑚𝑥 as proven earlier. In this case, 𝑦(𝑥) = 𝑥𝑒 𝑚𝑥 is proven as one of the
solution of 2nd order ODE and it is linearly independent with 𝑒 𝑚𝑥 . Thus a complete solution
𝑦(𝑥) = 𝑐1 𝑒 𝑚1 𝑥 + 𝑐2 𝑥𝑒 𝑚2 𝑥 is obtained which satisfy the linearly independency property.
In previous case, we discussed about homogeneous linear differential equation with constant coefficient, i.e.
𝑑2 𝑦 𝑑𝑦
𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 0 and let the solution to be 𝑦(𝑥) = 𝑒 𝑚𝑥 𝑜𝑟 𝑦(𝑥) = 𝑥𝑒 𝑚𝑥 depending on the roots of
characteristic equation.
𝑑2 𝑦
However this methods is inefficient to solve the Euler-Cauchy differential equation, i.e. 𝑥 2 +
𝑑𝑥 2
𝑑𝑦
𝑎𝑥 𝑑𝑥 + 𝑏𝑦 = 0 where the coefficients are not constant. The strategy to solve this type of differential
equation is to convert the non-constant coefficient into constant form. This can be achieved by
substitution (let 𝑥 = 𝑒 𝑡 ).
𝑑2 𝑦 𝑑2𝑦 𝑑𝑦
(i) 𝑥 2 𝑑𝑥 2 = 𝑑𝑡 2
− 𝑑𝑡
𝑑𝑦 𝑑𝑦
(ii) 𝑥 =
𝑑𝑥 𝑑𝑡
121
The detail description and proof is provided in the table below.
𝒅𝒚 𝒅𝒚
(i) Convert non-constant coefficient (𝒙 ) to constant coefficient ( )
𝒅𝒙 𝒅𝒕
𝑥 = 𝑒𝑡
>> 𝑙𝑛|𝑥| = 𝑡
1 𝑑𝑡
>> =
𝑥 𝑑𝑥
1 𝑑𝑡
>> 𝑥 = 𝑑𝑥
𝑑𝑦 𝑑𝑦 𝑑𝑡
Using chain rule, 𝑑𝑥
= 𝑑𝑡 𝑑𝑥
𝑑𝑦 𝑑𝑦 1
>> = ( )
𝑑𝑥 𝑑𝑡 𝑥
𝑑𝑦 𝑑𝑦
>> 𝑥 𝑑𝑥 = 𝑑𝑡
𝒅𝟐 𝒚 𝒅𝟐 𝒚 𝒅𝒚
(ii) Convert non-constant coefficient (𝒙𝟐 𝟐 ) to constant coefficient ( 𝟐 − )
𝒅𝒙 𝒅𝒕 𝒅𝒕
𝑑𝑦 𝑑𝑦
𝑥 𝑑𝑥 = 𝑑𝑡
𝑑 𝑑𝑦 𝑑 𝑑𝑦
>> 𝑑𝑥 (𝑥 𝑑𝑥 ) = 𝑑𝑥 ( 𝑑𝑡 )
𝑑2 𝑦 𝑑𝑦 𝑑 𝑑𝑦
>> 𝑥 𝑑𝑥 2 + 𝑑𝑥 = 𝑑𝑥 ( 𝑑𝑡 )
𝑑 𝑑𝑦 𝑑 𝑑𝑦 𝑑𝑡 1 𝑑𝑡
Using chain rule, 𝑑𝑥 ( 𝑑𝑡 ) = 𝑑𝑡 ( 𝑑𝑡 ) (𝑑𝑥) where 𝑥 = 𝑑𝑥
𝑑 𝑑𝑦 𝑑2 𝑦 1
>> 𝑑𝑥 ( 𝑑𝑡 ) = ( 𝑑𝑡 2 ) (𝑥)
𝑑2 𝑦 𝑑𝑦 𝑑2 𝑦 1
Combining the equations, we get 𝑥 𝑑𝑥 2 + 𝑑𝑥 = ( 𝑑𝑡 2 ) (𝑥)
𝑑2 𝑦 𝑑𝑦 𝑑2 𝑦 𝑑𝑦 𝑑𝑦
>> 𝑥 2 𝑑𝑥 2 + 𝑥 𝑑𝑥 = 𝑑𝑡 2
where 𝑥 𝑑𝑥 = 𝑑𝑡
𝑑2 𝑦 𝑑2 𝑦 𝑑𝑦
>> 𝑥 2 𝑑𝑥 2 = 𝑑𝑡 2
− 𝑑𝑡
122
𝑑2 𝑦 𝑑𝑦
For example: Solve 2𝑥 2 𝑑𝑥 2 − 3𝑥 𝑑𝑥 − 3𝑦 = 0.
𝑑2 𝑦 𝑑2𝑦 𝑑𝑦
(i) 𝑥2 2 = 2 −
𝑑𝑥 𝑑𝑡 𝑑𝑡
𝑑𝑦 𝑑𝑦
(ii) 𝑥 =
𝑑𝑥 𝑑𝑡
𝑑2𝑦 𝑑𝑦 𝑑2 𝑦 𝑑𝑦
2𝑥 2 𝑑𝑥 2 − 3𝑥 𝑑𝑥 − 3𝑦 = 0 [Euler-Cauchy differential equation, 𝑥 2 𝑑𝑥 2 + 𝑎𝑥 𝑑𝑥 + 𝑏𝑦 = 0]
𝑑2 𝑦 𝑑𝑦 𝑑𝑦
>> 2 ( 𝑑𝑡 2 − 𝑑𝑡 ) − 3 ( 𝑑𝑡 ) − 3𝑦 = 0
𝑑2 𝑦 𝑑𝑦
>> 2 ( 𝑑𝑡 2 ) − 5 ( 𝑑𝑡 ) − 3𝑦 = 0 [2nd order linear homogeneous DE with Constant coefficient]
where
𝑑2 𝑦 𝑑𝑦
Back substitution, we get the complementary solution to the 2𝑥 2 𝑑𝑥 2 − 3𝑥 𝑑𝑥 − 3𝑦 = 0.
>> 𝑥 = 𝑒 𝑡
123
12.4.3 NONHOMOGENEOUS LINEAR DIFFERENTIAL EQUATION WITH CONSTANT
COEFFICIENTS 𝑎, 𝑏, 𝑐 IN THE FORM OF 𝑟(𝑥) = 𝑒 𝛼𝑥 𝑃𝑛 (𝑥)
So far we have discussed two strategies to solve homogeneous problem, now we will continue with the
𝑑2 𝑦 𝑑𝑦
nonhomogeneous problem, i.e. 𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 𝑟(𝑥).
If the RHS components, 𝑟(𝑥) are in the form of exponential, polynomial, sine and cosine functions,
we can implement the method of undetermined coefficient by let the RHS components to be equal to
𝑒 𝛼𝑥 𝑃𝑛 (𝑥) as following:
𝑑2 𝑦 𝑑𝑦
𝑎 𝑑𝑥 2 + 𝑏 𝑑𝑥 + 𝑐𝑦 = 𝑒 𝛼𝑥 𝑃𝑛 (𝑥)
Hence, we can propose the possible particular solution of 𝑦𝑝 = 𝑒 𝛼𝑥 𝑄𝑛 (𝑥). The general procedure to
solve the 2nd order nonhomogeneous linear differential equation using the method of undetermined
coefficient is given below.
124
General procedure for the method of undetermined coefficient:
Step 5: The total solution for the 2nd order nonhomogeneous linear differential equation:
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝
Hint: Always solve the complementary solution first before proposing the particular solution.
The detail description for the “Step 1: Solve the homogenous part first” can be found in the previous section.
Now, we will discuss on the “Step 2: Solve the nonhomogenous part next”.
The method of undetermined coefficient is only applicable to 2nd order nonhomogeneous linear
differential equation, where the RHS components, 𝑟(𝑥) are restricted for exponential, polynomial, sine and
cosine functions, i.e. 𝑒 𝛼𝑥 𝑃𝑛 (𝑥).
125
The exponential function is related directly to the 𝑒 𝛼𝑥 and polynomial function is related directly to
the 𝑃𝑛 (𝑥). Moreover, the exponential function, 𝑒 𝛼𝑥 is related indirectly to sine & cosine functions through
Euler’s Formula: 𝑒 ±𝑖𝑥 = 𝑐𝑜𝑠𝑥 ± 𝑖(𝑠𝑖𝑛𝑥).
For example:
Depends on the RHS function, the possible particular solution is proposed for the 2nd order
nonhomogeneous linear differential equation as shown in table below.
126
(iv) Mixture of 𝑟(𝑥) 𝑦𝑝 = 𝑒 (6𝑖−3)𝑥 𝑄𝑂
(6𝑖)𝑥
Exponential & = 𝑅𝑒(𝑒 )(𝑒 (−3𝑥) )𝑃0 (𝑥) = 𝐴𝑒 (6𝑖−3)𝑥 (ii) 𝑟(𝑥) = 𝑐𝑜𝑠6𝑥
Cosine Function, = 𝑅𝑒(𝑒 (6𝑖−3)𝑥
)𝑃0 (𝑥) 𝑦𝑝
e.g. = 𝐶𝑐𝑜𝑠6𝑥 + 𝐷𝑠𝑖𝑛6𝑥
𝑟(𝑥) = 𝑒 −3𝑥 𝑐𝑜𝑠6𝑥 where 𝑦𝑝,𝑎𝑐𝑡𝑢𝑎𝑙 = 𝑅𝑒(𝑦𝑝 )
𝛼 = 6𝑖 − 3 (iii) 𝑟(𝑥) = 𝑒 −3𝑥 𝑐𝑜𝑠6𝑥
& 𝑦𝑝
Given 𝑃𝑛 (𝑥) = 1 with = 𝑒 −3𝑥 (𝐶𝑐𝑜𝑠6𝑥
+ 𝐷𝑠𝑖𝑛6𝑥)
𝑒 𝑖(6𝑥) = 𝑐𝑜𝑠6𝑥 + 𝑖(𝑠𝑖𝑛6𝑥) degree 𝑛 = 0
𝑅𝑒(𝑒 𝑖(6𝑥) ) = 𝑐𝑜𝑠6𝑥
(vi) Mixture of
𝑟(𝑥) = 𝑒 −3𝑥 𝑃1 (𝑥) 𝑦𝑝 = 𝑒 (−3)𝑥 𝑄1 Nil
Polynomial &
Exponential = 𝑒 (−3)𝑥 (𝐴𝑥 + 𝐵)
Function in where
multiplication, eg.
𝛼 = −3
𝑟(𝑥) = 6𝑥𝑒 −3𝑥 &
𝑃1 (𝑥) = 6𝑥
(vii) Mixture of
For polynomial function, 𝑦𝑝,1 = 𝑒 (0)𝑥 𝑄3 Can be solved
Polynomial & = 𝐴𝑥 3 + 𝐵𝑥 2 separately and then
Exponential 𝑟(𝑥) = 𝑒 (0𝑥) 𝑃3 (𝑥) +𝐶𝑥 + 𝐷 combine the result
Function in ‘+’ where 𝛼 = 0 & 𝑃3 (𝑥) = −3𝑥
later.
6𝑥 3 + 4𝑥 2 + 5 is the 𝑦𝑝,2 = 𝑒 𝑄𝑂
𝑟(𝑥) = 𝑒 −3𝑥 + 6𝑥 3 polynomial function of = 𝐸𝑒 −3𝑥
+4𝑥 2 + 5 degree 𝑛 = 3
For exponential function, 𝑦𝑝 = 𝑦𝑝,1 + 𝑦𝑝,2
𝑟(𝑥) = 𝑒 −3𝑥 𝑃0 (𝑥)
where 𝛼 = −3 & 𝑃𝑛 (𝑥) = 1
with degree 𝑛 = 0
Note: 𝑒 ±𝑖𝑥 = 𝑐𝑜𝑠𝑥 ± 𝑖(𝑠𝑖𝑛𝑥); 𝑄𝑛 (𝑥) & 𝑃𝑛 (𝑥) are two polynomial functions with same degree.
127
Now, we will discuss on the “Step 3: To check the linear dependency and give treatment to particular solution
if needed”. The possible particular solution is proposed according to the RHS function, however, further
treatment will be needed to obtain a linearly independent solution by comparing the complementary
solution. In fact, the actual particular solution 𝑦𝑝 can be separated into 3 cases depending on
Step 4 & 5 are quite straight forward, the solution of nonhomogeneous differential equation consists of
complementary solution and particular solution (i.e. 𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 ). Example of solving the 2nd order
nonhomogeneous linear using the method of undetermined coefficient is given below.
128
𝑑2 𝑦 𝑑𝑦
For example: Solve 𝑑𝑥 2
− 3 𝑑𝑥 + 2𝑦 = 𝑒 𝑥 [RHS - Pure Exponential Function]
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 2𝑥 + 𝑐2 𝑒 𝑥 − 𝑥𝑒 𝑥
129
𝑑2 𝑦 𝑑𝑦
For example: Solve 𝑑𝑥 2 − 5 𝑑𝑥 + 6𝑦 = 4𝑠𝑖𝑛2𝑥 [RHS - Pure Sine Function]
Step 1: Homogeneous Part Step 2: Nonhomogeneous Part
𝑑2 𝑦 𝑑𝑦 𝑑2 𝑦 𝑑𝑦
i.e. 𝑑𝑥 2
− 5 𝑑𝑥 + 6𝑦 = 0 i.e. 𝑑𝑥 2 − 5 𝑑𝑥 + 6𝑦 = 4𝑠𝑖𝑛2𝑥
130
The particular solution:
2
𝑦𝑝 = 1−5𝑖 𝑒 2𝑖𝑥
2 1+5𝑖
>> 𝑦𝑝 = 1−5𝑖 (1+5𝑖) (𝑐𝑜𝑠2𝑥 + 𝑖𝑠𝑖𝑛2𝑥)
2(1+5𝑖)
>> 𝑦𝑝 = (𝑐𝑜𝑠2𝑥 + 𝑖𝑠𝑖𝑛2𝑥)
26
(1+5𝑖)
>> 𝑦𝑝 = (𝑐𝑜𝑠2𝑥 + 𝑖𝑠𝑖𝑛2𝑥)
13
(𝑐𝑜𝑠2𝑥−5𝑠𝑖𝑛2𝑥)+𝑖(5𝑐𝑜𝑠2𝑥+𝑠𝑖𝑛2𝑥)
>> 𝑦𝑝 = 13
𝑑2 𝑦 𝑑𝑦
The complete/ general solution to −5 + 6𝑦 = 4𝑠𝑖𝑛2𝑥 is
𝑑𝑥 2 𝑑𝑥
(5𝑐𝑜𝑠2𝑥+𝑠𝑖𝑛2𝑥)
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 2𝑥 + 𝑐2 𝑒 3𝑥 + 13
131
𝑑2 𝑦 𝑑𝑦
For example: Solve 𝑑𝑥 2 + 4 𝑑𝑥 + 5𝑦 = 𝑒 −2𝑥 𝑐𝑜𝑠𝑥 [RHS - Mixture of Exponential & Cosine Function]
𝑑2 𝑦 𝑑𝑦
Substitute to the ODE equation: 𝑑𝑥 2 + 4 𝑑𝑥 + 5𝑦 = 𝑒 −2𝑥 𝑐𝑜𝑠𝑥
𝑑2 𝑦 𝑑𝑦
The complete/ general solution to 𝑑𝑥 2 + 4 𝑑𝑥 + 5𝑦 = 𝑒 −2𝑥 𝑐𝑜𝑠𝑥 is
1
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 (−2+𝑖)𝑥 + 𝑐2 𝑒 (−2−𝑖)𝑥 + 2 𝑥𝑒 −2𝑥 (𝑠𝑖𝑛𝑥)
There is another alternative to solve the 2nd order nonhomogeneous linear ODE problem with RHS sine and
cosine functions by using 𝑦𝑝 = 𝐶𝑐𝑜𝑠𝑥 + 𝐷𝑠𝑖𝑛𝑥 . Students are allowed to use any of them to solve the
problem.
𝑑2 𝑦 𝑑𝑦
2nd Alternative Method is used to solve the same example: Solve 𝑑𝑥 2 + 4 𝑑𝑥 + 5𝑦 = 𝑒 −2𝑥 𝑐𝑜𝑠𝑥
133
Complementary solution: Solve the coefficient for the particular solution:
𝑦𝑐 = 𝑐1 𝑒 (−2+i)𝑥 + 𝑐2 𝑒 (−2−i)𝑥 𝑦𝑝 = 𝑥𝑒 −2𝑥 (𝐶𝑐𝑜𝑠𝑥 + 𝐷𝑠𝑖𝑛𝑥)
Differentiate it, we get:
Using Euler equation:
𝑑𝑦𝑝
𝑦𝑐 = 𝑒 −2𝑥 (𝑐1 𝑒 𝑖𝑥 + 𝑐2 𝑒 −𝑖𝑥 ) = 𝑒 −2𝑥 (𝐶𝑐𝑜𝑠𝑥 + 𝐷𝑠𝑖𝑛𝑥) + 𝑥(−2𝑒 −2𝑥 )(𝐶𝑐𝑜𝑠𝑥 + 𝐷𝑠𝑖𝑛𝑥)
𝑑𝑥
Let 𝐴 = 𝑐1 + 𝑐2 & 𝐵 = 𝑐1 𝑖 + 𝑐2 𝑖 𝑑 2 𝑦𝑝
𝑑𝑥 2
= (−2𝑒 −2𝑥 )(𝐶𝑐𝑜𝑠𝑥 + 𝐷𝑠𝑖𝑛𝑥) + 𝑒 −2𝑥 (−𝐶𝑠𝑖𝑛𝑥 + 𝐷𝑐𝑜𝑠𝑥)
+𝑒 −2𝑥 ((−2𝐶 + 𝐷)𝑐𝑜𝑠𝑥 + (−2𝐷 + 𝐶)𝑠𝑖𝑛𝑥)
−2𝑥 (𝐴𝑐𝑜𝑠𝑥
𝑦𝑐 = 𝑒 + 𝐵𝑠𝑖𝑛𝑥) +𝑥(−2𝑒 −2𝑥 )((−2𝐶 + 𝐷)𝑐𝑜𝑠𝑥 + (−2𝐷 + 𝐶)𝑠𝑖𝑛𝑥)
+𝑥𝑒 −2𝑥 ((−2𝐶 + 𝐷)(−𝑠𝑖𝑛𝑥) + (−2𝐷 + 𝐶)𝑐𝑜𝑠𝑥)
= 𝑒 −2𝑥 ((−4𝐶 + 2𝐷)𝑐𝑜𝑠𝑥 + (−4𝐷 − 2𝐶)𝑠𝑖𝑛𝑥)
+𝑥𝑒 −2𝑥 ((3𝐶 − 4𝐷)𝑐𝑜𝑠𝑥 + (3𝐷 + 4𝐶)𝑠𝑖𝑛𝑥)
𝑑2 𝑦 𝑑𝑦
Substitute to the ODE equation: 𝑑𝑥 2 + 4 𝑑𝑥 + 5𝑦 = 𝑒 −2𝑥 𝑐𝑜𝑠𝑥
𝑑2 𝑦 𝑑𝑦
The complete/ general solution to 𝑑𝑥 2 + 4 𝑑𝑥 + 5𝑦 = 𝑒 −2𝑥 𝑐𝑜𝑠𝑥 is
1
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑒 −2𝑥 (𝐴𝑐𝑜𝑠𝑥 + 𝐵𝑠𝑖𝑛𝑥) + 2 𝑥𝑒 −2𝑥 (𝑠𝑖𝑛𝑥)
134
𝑑2 𝑦
For example: Solve 𝑑𝑥 2 + 4𝑦 = 8𝑥 2 [RHS - Pure Polynomial Function]
Comment:
Complementary solution:
(i) 𝑦𝑝 = 𝐴2 𝑥 2 + 𝐴1 𝑥 + 𝐴0 and 𝑦𝑐 = 𝑐1 𝑒 (2i)𝑥 + 𝑐2 𝑒 (−2i)𝑥 are
(2i)𝑥 (−2i)𝑥
𝑦𝑐 = 𝑐1 𝑒 + 𝑐2 𝑒 linearly independent
𝑑2 𝑦
The complete/ general solution to 𝑑𝑥 2 + 4𝑦 = 8𝑥 2 is
Comment:
Complementary solution:
(i) 𝑦𝑝 = (𝐴𝑥 + 𝐵)𝑒 4𝑥 and 𝑦𝑐 = 𝑐1 𝑒 6𝑥 + 𝑐2 𝑒 −2𝑥 are linearly
6𝑥 −2𝑥
𝑦𝑐 = 𝑐1 𝑒 + 𝑐2 𝑒 independent
Solve the coefficient for the particular solution:
𝑦𝑝 = (𝐴𝑥 + 𝐵)𝑒 4𝑥
Differentiate it, we get:
𝑑𝑦𝑝
= 4(𝐴𝑥 + 𝐵)𝑒 4𝑥 + (𝐴)𝑒 4𝑥
𝑑𝑥
𝑑 2 𝑦𝑝
𝑑𝑥 2
= 16(𝐴𝑥 + 𝐵)𝑒 4𝑥 + 8(𝐴)𝑒 4𝑥
𝑑2 𝑦 𝑑𝑦
Substitute to the ODE equation: 𝑑𝑥 2 − 4 𝑑𝑥 − 12𝑦 = 𝑥𝑒 4𝑥
𝑑2 𝑦 𝑑𝑦
The complete/ general solution to 𝑑𝑥 2 − 4 𝑑𝑥 − 12𝑦 = 𝑥𝑒 4𝑥 is
1 1
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 6𝑥 + 𝑐2 𝑒 −2𝑥 + (−12 𝑥 − 36) 𝑒 4𝑥
136
𝑑2 𝑦 𝑑𝑦
For example: Solve 𝑑𝑥 2 − 4 𝑑𝑥 + 3𝑦 = 3𝑒 2𝑥 + 4𝑥 [RHS - Mixture of Polynomial & Exponential Function in ‘+’]
137
𝑑2 𝑦 𝑑𝑦
The complete/ general solution to 𝑑𝑥 2
− 4 𝑑𝑥 + 3𝑦 = 3𝑒 2𝑥 + 4𝑥 is
4 16
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 𝑥 + 𝑐2 𝑒 3𝑥 − 3𝑒 2𝑥 + 3 𝑥 + 9
Hint: The example above illustrates the linearity or superposition principle, where the solutions can be
added directly as illustrated below.
𝑑2 𝑦 𝑑𝑦
>> 𝑦𝑝,1 = −3𝑒 2𝑥 is the particular solution to − 4 + 3𝑦 = 3𝑒 2𝑥 ;
𝑑𝑥 2 𝑑𝑥
4 16 𝑑2 𝑦 𝑑𝑦
>> 𝑦𝑝,2 = 3 𝑥 + 9 is the particular solution to 𝑑𝑥 2 − 4 𝑑𝑥 + 3𝑦 = 4𝑥;
𝑑2 𝑦 𝑑𝑦
>> 𝑦𝑝,𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑝,1 + 𝑦𝑝,2 is the total particular solution to 2 − 4 + 3𝑦 = 3𝑒 2𝑥 + 4𝑥.
𝑑𝑥 𝑑𝑥
For example: A 𝑅𝐿𝐶 Circuit is provided as shown in the figure below where there are an inductor of 𝐿 =
50 ℎ𝑒𝑛𝑟𝑦𝑠, a resistor of 𝑅 = 5 𝑜ℎ𝑚𝑠 and a capacitor of 𝐶 = 8 𝑓𝑎𝑟𝑎𝑑𝑠. At 𝑡 = 0, the switch is closed. Given
𝑑 2 𝑞(𝑡) 𝑑𝑞(𝑡) 1
the 2nd order ODE for the system is 𝐿 +𝑅 + 𝑞(𝑡) = 𝐸(𝑡).
𝑑𝑡 2 𝑑𝑡 𝐶
Find the charge and current at any time 𝑡 > 0 if (a) The voltage is supplied by a DC battery, i.e. 𝐸 = 40 𝑣𝑜𝑙𝑡𝑠
138
(a) The voltage is supplied by a DC battery, i.e. 𝐸 = 40 𝑣𝑜𝑙𝑡𝑠
50𝑚2 + 5𝑚 + 8 = 0
1 RHS : 𝑟(𝑥) = 𝑒 𝛼𝑥 𝑃𝑛 (𝑥) = 40
where 𝛼 = 0, 𝑛 = 0
400𝑚2 + 40𝑚 + 1 = 0
(20𝑚 + 1)2 = 0 Possible particular solution:
>>𝐴 = 320
The actual particular solution:
𝑞𝑝 = 320
𝑑 2 𝑞(𝑡) 𝑑𝑞(𝑡) 1
The complete/ general solution to 50 𝑑𝑡 2 + 5 𝑑𝑡 + 8 𝑞(𝑡) = 40 is
(i) Charge solution,
𝑞𝑡𝑜𝑡𝑎𝑙 = 𝑞𝑐 + 𝑞𝑝 = 𝑐1 𝑒 −0.05𝑡 + 𝑐2 𝑡𝑒 −0.05𝑡 + 320
139
Solution to Initial value problem
Note:
The general solution is not the actual solution to the problem because it has infinite 𝑐1 , 𝑐2 that can
satisfy the problem. Recall that in the initial value problem, we can further solve the constant 𝑐1 , 𝑐2
in the system if the initial condition of the problem is known in advance.
Note: We plot the result of charge and current solutions (for the 𝑅𝐿𝐶 circuit problem with 𝐿 = 50 ℎ𝑒𝑛𝑟𝑦𝑠,
a resistor of 𝑅 = 5 𝑜ℎ𝑚𝑠 and a capacitor of 𝐶 = 8 𝑓𝑎𝑟𝑎𝑑𝑠) as follows. The intention of this example is to
encourage student to link the mathematical result to the actual problem instead of calculating it for nothing.
However, data analysis on specific engineering problem requires certain knowledge on the subject and
hence it is out of the scope in this study.
140
Hint: Relate this to the scenario of charging battery.
(1) Identify the transient and steady state region in the graph. Why it is important to identify them?
Ans: Transient region happens within around 0 − 150𝑠; while steady state region occurs after that.
This relationship is important to the RLC circuit problem such as charging battery, where we can
estimate how much time is needed to fully charge the battery. Moreover, it shows that the current is
reduced to zero once it is fully charged.
(2) What are the relationships between complementary solution & particular solution with these
transient and steady state region?
Ans: We know that the total solution is a combination of the complementary solution & particular
solution. The complementary solution contributes more within the transient region and its effect
diminishes within the steady state region; while the particular solution contributes significantly within
the steady state region.
(3) Why does the charge behave as such: increase over a time initially and at time approximately 150s,
the charge remains constant afterward? Why does the current behave as such: increase over a time
initially and decrease after it reaches its maximum? The charge decreases to zero at time approximately
150s and remains constant afterward.
Ans: To understand the changes, we need to check the equation, 𝑞𝑡𝑜𝑡𝑎𝑙 = −320𝑒 −0.05𝑡 −
16𝑡𝑒 −0.05𝑡 + 320, where complementary solution of charge, 𝑞𝑐 = −320𝑒 −0.05𝑡 − 16𝑡𝑒 −0.05𝑡 and its
particular solution, 𝑞𝑝 = 320.Since the complementary solution consists the exponential function,
i.e. 𝑒 −0.05𝑡 . As the time increase, this function will approach zero and thus diminish. At the same time,
the particular solution remains all the time. Hence, 𝑞𝑡𝑜𝑡𝑎𝑙 = −320𝑒 −0.05𝑡 − 16𝑡𝑒 −0.05𝑡 + 320
increase over a time initially and at time approximately 150s, the charge remains constant afterward.
𝑑𝑞𝑡𝑜𝑡𝑎𝑙
The current is the rate of change of the charge, i.e. 𝑖𝑡𝑜𝑡𝑎𝑙 = 𝑑𝑡
= 𝑒 −0.05𝑡 (16𝑡). Since it consists of
the exponential function, it illustrates that the charging rate will be higher initially and reduce to zero
afterwards, as this function will approach to zero as the time increase.
Think: In case you want to reduce the charging time or increase the charging capacity, what should you do?
141
12.4.4 SOLVING ALL TYPES OF NONHOMOGENEOUS LINEAR DIFFERENTIAL EQUATION
WITH METHOD OF VARIATION OF PARAMETERS
The previous undetermined coefficients method is only applicable for simple functions such as a mixture of
exponential and polynomial functions. However, it is impractical to solve complicated function other than
exponential and polynomial functions such as tangent function, Mixture of Polynomial & Exponential
Function in ‘÷’, logarithmic function, etc. To solve 2nd order ODE with complicated function in its RHS, we
recommend user to use the method of variation of parameters.
𝑑2 𝑦 𝑑𝑦
(1) Standard form: 𝑑𝑥 2
+ 𝑝(𝑥) 𝑑𝑥 + 𝑞(𝑥)𝑦 = 𝑟(𝑥)
(2) Compute the complementary solution, 𝑦𝑐 = 𝑐1 𝑦1 (𝑥) + 𝑐2 𝑦2 (𝑥) by using method introduced in section
12.4.1 & 12.4.2, where 𝑐1 & 𝑐2 are arbitrary constants.
The drawback of this technique is that it is time consuming to complete the integration operation
and there are cases where the integration function cannot be solved analytically (using calculus). In this case,
we may need an advance tool such as numerical method to solve the integration problem.
142
𝑑2 𝑦 𝑑𝑦
For example: Solve 𝑑𝑥 2 + 2 𝑑𝑥 + 𝑦 = 𝑒 −𝑥 𝑙𝑛𝑥
𝑦𝑐 = 𝑐1 𝑒 −𝑥 + 𝑐2 𝑥𝑒 −𝑥 (𝑒 −𝑥 𝑙𝑛𝑥)(𝑥𝑒 −𝑥 )
= −∫ 𝑒 −2𝑥
𝑑𝑥
= − ∫ 𝑥𝑙𝑛𝑥 𝑑𝑥
Comment: Using integration by part method:
(i) 𝑦𝑐,1 = 𝑒 −𝑥 and 𝑦𝑐,2 = 𝑒 −𝑥 are linearly Let 𝑢 = 𝑙𝑛𝑥; 𝑑𝑣 = 𝑥𝑑𝑥
dependent.
1 𝑥2
(ii) Treatment is done so that 𝑦𝑐,1 = 𝑑𝑢 = 𝑥 𝑑𝑥; 𝑣 = 2
𝑒 −𝑥 and 𝑦𝑐,2 = 𝑥𝑒 −𝑥 are linearly
𝑢1 (𝑥) = − ∫ 𝑥𝑙𝑛𝑥 𝑑𝑥 = −(𝑢𝑣 − ∫ 𝑣𝑑𝑢)
independent.
𝑥2 𝑥2 1
= − (𝑙𝑛𝑥 ( 2 ) − ∫ 2 𝑥
𝑑𝑥)
Compute the Wronskian, 𝑊(𝑦1 , 𝑦2 )
𝑦1 𝑦2 𝑥2 𝑥
= − (𝑙𝑛𝑥 ( 2 ) − ∫ 2 𝑑𝑥)
= |𝑑𝑦1 𝑑𝑦2 |
𝑑𝑥−𝑥 𝑑𝑥 = (−𝑙𝑛𝑥 ( ) +
𝑥2 𝑥2
)
𝑒 𝑥𝑒 −𝑥 2 4
= | −𝑥 |
−𝑒 −𝑥𝑒 −𝑥 + 𝑒 −𝑥 𝑟(𝑥)𝑦1 (𝑥)
=𝑒 −𝑥 (−𝑥𝑒 −𝑥
+𝑒 −𝑥 )
− (𝑥𝑒 −𝑥 )(−𝑒 −𝑥 ) 𝑢2 (𝑥) = ∫ 𝑑𝑥
𝑊(𝑦1 ,𝑦2 )
−2𝑥
=𝑒
(𝑒 −𝑥 𝑙𝑛𝑥)(𝑒 −𝑥 )
=∫ 𝑒 −2𝑥
𝑑𝑥
= ∫ 𝑙𝑛𝑥 𝑑𝑥
Using integration by part method:
1
Let 𝑢 = 𝑙𝑛𝑥; 𝑑𝑣 = 𝑑𝑥, then 𝑑𝑢 = 𝑥 𝑑𝑥; 𝑣 = 𝑥
= (𝑙𝑛𝑥(𝑥) − ∫ 1𝑑𝑥)
= (𝑙𝑛𝑥(𝑥) − 𝑥)
143
The particular solution:
𝑦𝑝 = 𝑢1 (𝑥)𝑦1 (𝑥) + 𝑢2 (𝑥)𝑦2 (𝑥)
𝑥2 𝑥2
= (−𝑙𝑛𝑥 ( 2 ) + 4
) 𝑒 −𝑥 + (𝑙𝑛𝑥(𝑥) − 𝑥)(𝑥𝑒 −𝑥 )
𝑥2 3𝑥 2
= (𝑙𝑛𝑥 ( 2 ) − 4
) 𝑒 −𝑥
𝑑2 𝑦 𝑑𝑦
The complete/ general solution to 𝑑𝑥 2 + 2 𝑑𝑥 + 𝑦 = 𝑒 −𝑥 𝑙𝑛𝑥 is
𝑥2 3𝑥 2
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 −𝑥 + 𝑐2 𝑥𝑒 −𝑥 + (𝑙𝑛𝑥 ( 2 ) − 4
) 𝑒 −𝑥
144
𝑑2 𝑦 𝑑𝑦 𝑒𝑥
For example: Solve 𝑑𝑥 2 − 2 𝑑𝑥 + 𝑦 = 𝑥 2 +1
145
The particular solution:
𝑦𝑝 = 𝑢1 (𝑥)𝑦1 (𝑥) + 𝑢2 (𝑥)𝑦2 (𝑥)
1
= (− 2 𝑙𝑛|𝑥 2 + 1|) 𝑒 𝑥 + (𝑡𝑎𝑛−1 𝑥)(𝑥𝑒 𝑥 )
1
= (𝑥𝑡𝑎𝑛−1 𝑥 − 2 𝑙𝑛|𝑥 2 + 1|) 𝑒 𝑥
𝑑2 𝑦 𝑑𝑦 𝑒𝑥
The complete/ general solution to 𝑑𝑥 2 − 2 𝑑𝑥 + 𝑦 = 𝑥 2 +1 is
1
𝑦𝑡𝑜𝑡𝑎𝑙 = 𝑦𝑐 + 𝑦𝑝 = 𝑐1 𝑒 𝑥 + 𝑐2 𝑥𝑒 𝑥 + (𝑥𝑡𝑎𝑛−1 𝑥 − 𝑙𝑛|𝑥 2 + 1|) 𝑒 𝑥
2
146
POWER SERIES SOLUTIONS FOR
DIFFERENTIAL EQUATIONS
WEEK 13: POWER SERIES SOLUTIONS FOR DIFFERENTIAL EQUATIONS
13.1 POWER SERIES METHOD
Power Series
A power series is an infinite series of the form
∞
∑ 𝑎𝑛 (𝑥 − 𝑥0 )𝑛
𝑛=0
= 𝑎0 + 𝑎1 (𝑥 − 𝑥0 ) + 𝑎2 (𝑥 − 𝑥0 )2 + ⋯ (1)
where a0, a1, a2, ... are real constants, called the coefficients of the series, x0 is a constant, called the center
of the series, and x is a variable.
∑ 𝑎𝑛 𝑥 𝑛 = 𝑎0 + 𝑎1 𝑥 + 𝑎2 𝑥 2 + 𝑎3 𝑥 3 + ⋯
𝑛=0
∞
𝑥
𝑥𝑛 𝑥2 𝑥3
(ii) 𝑒 = ∑ = 1 + 𝑥 + + +⋯
𝑛! 2! 3!
𝑛=0
∞
(−1)𝑛 𝑥 2𝑛 𝑥2 𝑥4
(iii) cos 𝑥 = ∑ = 1 − + ±⋯
(2𝑛)! 2! 4!
𝑛=0
∞
(−1)𝑛 𝑥 2𝑛+1 𝑥3 𝑥5
(iv) sin 𝑥 = ∑ = 𝑥 − + ±⋯
(2𝑛 + 1)! 3! 5!
𝑛=0
147
∞
(−1)𝑛+1 𝑥 𝑛 𝑥2 𝑥3
(v) ln(1 + 𝑥) = ∑ =𝑥 − + − ⋯
𝑛 2 3
𝑛=1
Example:
For the geometric series
1 + 𝑥 + 𝑥2 + ⋯ + 𝑥𝑛 + ⋯
Then:
𝑠1 = 1 + 𝑥 𝑅1 = 𝑥 2 + 𝑥 3 + 𝑥 4 + ⋯
𝑠2 = 1 + 𝑥 + 𝑥 2 𝑅2 = 𝑥 2 + 𝑥 3 + 𝑥 4 + ⋯
etc.
If for some x = x1, sn(x) converges, that is, lim 𝑠𝑛 (𝑥1 ) = 𝑠(𝑥1 ) then the series (1) converges, or is called
𝑛→∞
convergent at x = x1; and the number s(x1) is called the value or sum of (1) at x1, and can be written as
∞
𝑠(𝑥1 ) = ∑ 𝑎𝑛 (𝑥1 − 𝑥0 )𝑛
𝑛=0
If the sequence is divergent at x = x1, then the series (1) is said to diverge, or to be divergent at x = x1.
Note:
1. The series (1) converges at x = x0 when all its terms except for the first a0 are zero. In unusual
cases this may be the only x for which (1) converges.
2. If there are further values of x for which the series (1) converges, these value form an interval,
called the convergence interval. If this interval is finite, it has the midpoint x0 so that it is of the
form
|𝑥 − 𝑥0 | < 𝑅
148
and the series (1) converges for all x such that |𝑥 − 𝑥0 | < 𝑅 and diverges for all x such that
|𝑥 − 𝑥0 | > 𝑅. The number R is called the radius of convergence of (1). It can be obtained from
either of the following formulas:
1 𝑎𝑛+1
(𝑎) 𝑅 = (𝑏) 𝑅 = lim | | (3)
𝑛
lim √|𝑎𝑚 | 𝑛→∞ 𝑎𝑛
𝑛→∞
provided these limits exist and are not zero. [If they are infinite, then (1) converges only at the
center x0.]
3. The convergence interval may sometimes be infinite, that is, (1) converges for all x. For example,
if the limit in (3a) and (3b) is zero. Then 𝑅 = ∞, for convenience.
4. Since power series are functions of x and we know that not every series will in fact exist, it then
makes sense to ask if a power series will exist for all x. This question is answered by looking at the
convergence of the power series. We say that a power series converges for x = c if the series,
∞
∑ 𝑎𝑛 (𝑐 − 𝑥0 )𝑛
𝑛=0
converges. Recall that this series will converge if the limit of partial sums,
𝑁
lim ∑ 𝑎𝑛 (𝑐 − 𝑥0 )𝑛
𝑁→∞
𝑛=0
exists and is finite. In other words, a power series will converge for x = c if
∞
∑ 𝑎𝑛 (𝑐 − 𝑥0 )𝑛
𝑛=0
is a finite number.
5. A power series will always converge if x = x0. In this case the power series will become
∞
∑ 𝑎𝑛 (𝑐 − 𝑥0 )𝑛 = 𝑎0
𝑛=0
With this it is known now that power series are guaranteed to exist for at least one value of x. The
following fact about the convergence of a power series is derived.
Fact
Given a power series, (1), there will exist a number 0 ≤ 𝜌 ≤ ∞ so that the power series will converge
for|𝑥 − 𝑥0 | < 𝜌 and diverge for |𝑥 − 𝑥0 | > 𝜌. This number is called the radius of convergence.
149
13.1.2 TEST FOR CONVERGENCE
Problem Set 1
Find the radius of convergence of the following series.
∞
𝑘
1. ∑
2𝑘
𝑘=1
Solution:
150
𝑘 + 1 2𝑘 𝑘+1 1 𝑘+1 1 1 1
𝜌 = lim | 𝑘+1 ∙ | = lim | | = lim | |= lim |1 + | =
𝑘→∞ 2 𝑘 𝑘→∞ 2𝑘 2 𝑘→∞ 𝑘 2 𝑘→∞ 𝑘 2
The series converges
∞
(−1)𝑘 (𝑥 − 3)𝑘
2. ∑
3𝑘 (𝑘 + 1)
𝑘=1
Solution:
(−1)𝑘+1 (𝑥 − 3)𝑘+1 3𝑘 (𝑘 + 1) (−1)(𝑥 − 3)(𝑘 + 1)
𝜌 = lim | 𝑘+1 ∙ | = lim | |
𝑘→∞ 3 (𝑘 + 1 + 1) (−1)𝑘 (𝑥 − 3)𝑘 𝑘→∞ 3(𝑘 + 2)
𝑦(𝑥) = ∑ 𝑎𝑛 (𝑥 − 𝑥0 )𝑛
𝑛=0
converges for |𝑥 − 𝑥0 | < 𝑅 where R > 0, then the series obtained by differentiating term by term
also converges for those x and represents the derivative y’ of y for those x, that is,
∞
151
Similarly,
∞
and so on.
have positive radii of convergence and their sums are f(x) and g(x), respectively, then the series
∞
∑(𝑎𝑛 + 𝑏𝑛 )(𝑥 − 𝑥0 )𝑛
𝑛=0
converges and represent f(x) + g(x) for each x that lies in the interior of the convergence interval
of each of the given series.
(3) Termwise multiplication
Two power series may be multiplied term by term. Suppose that
∞ ∞
𝑛
∑ 𝑎𝑛 (𝑥 − 𝑥0 ) and ∑ 𝑏𝑛 (𝑥 − 𝑥0 )𝑛
𝑛=0 𝑛=0
have positive radii of convergence and let f(x) and g(x) be their sums, respectively. Then the series
obtained by multiplying each term of the first series by each term of the second series and
collecting like powers of x – x0, that is,
∞
If a power series has a positive radius of convergence and a sum that is identically zero throughout its interval
of convergence, then each coefficient of the series is zero.
152
Sifting summation indices
(1) An index of summation is a dummy and can be changed.
Example:
∞ ∞
3𝑛 𝑛 2 3𝑘 𝑘 2 81
∑ =∑ = 1 + 18 + + ⋯.
𝑛! 𝑘! 2
𝑛=1 𝑘=1
as a single series; firstly, take x2 and 2, respectively, inside the summation, obtaining
∞ ∞
𝑛=2 𝑛=1
𝑠=2 𝑠=0
∑[𝑠(𝑠 − 1)𝑎𝑠 + 2(𝑠 + 1)𝑎𝑠+1 ] 𝑥 𝑠 = 2𝑎1 + 4𝑎2 𝑥 + (2𝑎2 + 6𝑎3 )𝑥 2 + (6𝑎3 + 8𝑎4 )𝑥 3 + ⋯
𝑠=0
153
Problem Set 2
Find the series of the following functions.
2
1. 𝑒𝑥
Solution:
∞
𝑥2
(𝑥 2 )𝑚 (𝑥 2 )2 (𝑥 2 )3
𝑒 = ∑ = 1 + 𝑥2 + + +⋯
𝑚! 2! 3!
𝑚=0
𝑥4 𝑥6
= 1 + 𝑥2 + + +⋯
2 6
2. 𝑒 𝑥 + sin 𝑥
Solution:
∞ ∞
𝑥
𝑥𝑚 (−1)𝑚 𝑥 2𝑚+1
𝑒 + sin 𝑥 = ∑ + ∑
𝑚! (2𝑚 + 1)!
𝑚=0 𝑚=0
𝑥2 𝑥3 𝑥3 𝑥5
=1 + 𝑥 + + + ⋯+ 𝑥 − + ±⋯
2! 3! 3! 5!
𝑥2 𝑥4 2𝑥 5
= 1 + 2𝑥 + + + +⋯
2 4! 5!
3. 𝑒 𝑥 (cos 𝑥)
Solution:
∞ ∞
𝑥
𝑥𝑚 (−1)𝑚 𝑥 2𝑚
𝑒 (cos 𝑥) = ( ∑ )(∑ )
𝑚! (2𝑚)!
𝑚=0 𝑚=0
2
𝑥 𝑥3 𝑥2 𝑥4
= (1 + 𝑥 + + + ⋯ ) (1 − + ± ⋯)
2! 3! 2! 4!
𝑥2 𝑥4 𝑥3 𝑥5
=1 − + + …+ 𝑥 − + +⋯
2! 4! 2! 4!
Before finding series solutions to differential equations; we need to determine when we can find series
solutions to differential equations with nonconstant coefficients. So, let’s start with the differential
equation,
𝑝(𝑥)𝑦 ′′ + 𝑞(𝑥)𝑦 ′ + 𝑟(𝑥)𝑦 = 0 (5)
154
To this point we’ve only dealt with constant coefficients. However, with series solutions we can now have
nonconstant coefficient differential equations. Also, here we will be dealing only with polynomial
coefficients.
Now, we say that x = x0 is an ordinary point if provided both
𝑞(𝑥) 𝑟(𝑥)
and
𝑝(𝑥) 𝑝(𝑥)
are analytic at x = x0. That is to say that these two quantities have Taylor series around x = x0. Since, we
are only dealing with coefficients that are polynomials so this will be equivalent to saying that
𝑝(𝑥0 ) ≠ 0
for most of the problems.
If a point is not an ordinary point we call it a singular point.
The basic idea to finding a series solution to a differential equation is to assume that we can write the
solution as a power series in the form,
∞
𝑦(𝑥) = ∑ 𝑎𝑛 (𝑥 − 𝑥0 )𝑛 (6)
𝑛=0
and then try to determine what the an’s need to be. We will only be able to do this if the point x = x0, is an
ordinary point. We will usually say that (6) is a series solution around x = x0.
Problem Set 3
𝑦 = ∑ 𝑎𝑛 (𝑥 − 0) = ∑ 𝑎𝑛 𝑥 𝑛
𝑛
𝑛=0 𝑛=0
Then,
∞
𝑦 = ∑ 𝑛 𝑎𝑛 𝑥 𝑛−1
′
𝑛=1
∞
𝑦 = ∑ 𝑛 (𝑛 − 1)𝑎𝑛 𝑥 𝑛−2
′′
𝑛=2
155
∞ ∞
𝑛−2
∑ 𝑛 (𝑛 − 1)𝑎𝑛 𝑥 − 𝑥 ∑ 𝑎𝑛 𝑥 𝑛 = 0
𝑛=2 𝑛=0
Step 3: Shift the first series down by 2 and the second series up by 1 to get both of the series in
terms of xn
∞ ∞
𝑛=0 𝑛=1
Step 4: Get the two series starting at the same value of n. The only way to do that for this problem is to
strip out the n = 0 term
∞ ∞
(2)(1)𝑎2 𝑥 + ∑(𝑛 + 2) (𝑛 + 1)𝑎𝑛+2 𝑥 − ∑ 𝑎𝑛−1 𝑥 𝑛 = 0
0 𝑛
𝑛=1 𝑛=1
∞
Step 5: Set all the coefficients equal to zero. The n = 0 coefficient is in front of the series and the n
= 1,2,3… are all in the series. So, setting coefficient equal to zero gives,
𝑛=0 2𝑎2 = 0
𝑛 = 1, 2, 3, … (𝑛 + 2)(𝑛 + 1)𝑎𝑛+2 − 𝑎𝑛−1 = 0
Step 6: Solving the first as well as the recurrence relation gives
𝑛=0 𝑎2 = 0
𝑎𝑛−1
𝑛 = 1, 2, 3, … 𝑎𝑛+2 −=
(𝑛 + 2)(𝑛 + 1)
Step 7: Start plugging in values of n
𝑎0
𝑛=1 𝑎3 =
(3)(2)
𝑎1
𝑛=2 𝑎4 =
(4)(3)
𝑎2
𝑛=3 𝑎5 = =0
(5)(4)
𝑎3 𝑎0
𝑛=4 𝑎6 = =
(6)(5) (6)(5)(3)(2)
𝑎4 𝑎1
𝑛=5 𝑎7 = =
(7)(6) (7)(6)(4)(3)
156
𝑎5
𝑛=6 𝑎8 = =0
(8)(7)
𝑎0
𝑎3𝑘 = 𝑘 = 1, 2, 3, ⋯
(2)(3)(5)(6) ⋯ (3𝑘 − 1)(3𝑘)
𝑎1
𝑎3𝑘+1 = 𝑘 = 1, 2, 3, ⋯
(3)(4)(6)(7) ⋯ (3𝑘)(3𝑘 + 1)
𝑎3𝑘+2 = 0 𝑘 = 0, 1, 2, ⋯
Note: Every third coefficient is zero. The formulas here are somewhat unpleasant and not all that easy
to see the first time around. These formulas will not work for k = 0.
Note: The series could not start at k = 0 since the general term doesn’t hold for k = 0
2. Find the first four terms in each portion of the series solution around x0 = -2 for the following
differential equation
𝑦 ′′ − 𝑥𝑦 = 0
Solution:
In this case, 𝑝(𝑥) = 1; hence for this differential equation every point is an ordinary point.
Assume solution:
∞ ∞
𝑦 = ∑ 𝑎𝑛 (𝑥 − (−2)) = ∑ 𝑎𝑛 (𝑥 + 2)𝑛
𝑛
𝑛=0 𝑛=0
Then,
∞
𝑦 = ∑ 𝑛 𝑎𝑛 (𝑥 + 2)𝑛−1
′
𝑛=1
157
∞
𝑦 = ∑ 𝑛 (𝑛 − 1)𝑎𝑛 (𝑥 + 2)𝑛−2
′′
𝑛=2
Step 2: Get all the coefficients moved into the series. There is a difference between this example and
the previous example. In this case we can’t just multiply the x into the second series since in
order to combine with the series it must be x + 2. Therefore we will first need to modify the
coefficient of the second series before multiplying it into the series.
∞ ∞
𝑛−2
∑ 𝑛 (𝑛 − 1)𝑎𝑛 (𝑥 + 2) − (𝑥 + 2 − 2) ∑ 𝑎𝑛 (𝑥 + 2)𝑛 = 0
𝑛=2 𝑛=0
∞ ∞ ∞
𝑛−2
∑ 𝑛 (𝑛 − 1)𝑎𝑛 (𝑥 + 2) − (𝑥 + 2) ∑ 𝑎𝑛 (𝑥 + 2) + 2 ∑ 𝑎𝑛 (𝑥 + 2)𝑛 = 0
𝑛
Step 3: Need to shift the first series down by 2 and the second series up by 1 to get common exponents
in all the series
∞ ∞ ∞
Step 4: Combine the series by stripping out the n = 0 terms from both the first and third series
∞ ∞ ∞
158
Step 6: Solve the first as well as the recurrence relation. In the first case there are two options, we can
solve for a2 or we can solve for a0. Out of habit I’ll solve for a0. In the recurrence relation we’ll
solve for the term with the largest subscript
𝑛=0 𝑎2 = −𝑎0
𝑎𝑛−1 − 2𝑎𝑛
𝑛 = 1, 2, 3, … 𝑎𝑛+2 =
(𝑛 + 2)(𝑛 + 1)
Note 1: This example we won’t be having every third term drop out as we did in the previous example.
Note 2: At this point we’ll also acknowledge that the instructions for this problem are different as well.
We aren’t going to get a general formula for the an’s this time so we’ll have to be satisfied with
just getting the first couple of terms for each portion of the solution. This is often the case for
series solutions. Getting general formulas for the an’s is the exception rather than the rule in
these kinds of problems.
Step 7: Start plugging in values of n. To get the first four terms we’ll just start plugging in terms until
we’ve got the required number of terms. Note that we will already be starting with an a0 and
an a1 from the first two terms of the solution so all we will need are three more terms with an a0
in them and three more terms with an a1 in them
𝑛=0 𝑎2 = −𝑎0
𝑎0 − 2𝑎1 𝑎0 𝑎1
𝑛=1 𝑎3 = = −
(3)(2) 6 3
𝑎1 − 2𝑎2 𝑎1 − 2(−𝑎0 ) 𝑎0 𝑎1
𝑛=2 𝑎4 = = = +
(4)(3) (4)(3) 6 12
𝑎2 − 2𝑎3 𝑎0 1 𝑎0 𝑎1 𝑎0 𝑎1
𝑛=3 𝑎5 = = − ( − )=− +
(5)(4) 20 10 6 3 15 30
Step 9: Collect up the terms that contain the same coefficient, factor the coefficient out and write the
results as a new series
1 1 1
𝑦(𝑥) = 𝑎0 {1 − (𝑥 + 2)2 + (𝑥 + 2)3 + (𝑥 + 2)4 − (𝑥 + 2)5 + ⋯ }
6 6 15
1 1 1
+ 𝑎1 {(𝑥 + 2) − (𝑥 + 2) + (𝑥 + 2)4 + (𝑥 + 2)5 + ⋯ }
3
3 12 30
159
Note: That’s the solution for this problem as far as we’re concerned. Notice that this solution looks
nothing like the solution to the previous example. It’s the same differential equation, but
changing x0 completely changed the solution.
3. Determine a series solution about x0 = 0 for the following initial value problem.
𝑦 ′′ − 2𝑥𝑦 ′ + 𝑦 = 0, 𝑦(0) = 1, 𝑦 ′ (0) = 1
Solution:
Assume solution:
∞
𝑦 = ∑ 𝑎𝑛 𝑥 𝑛
𝑛=0
Then,
∞
𝑦 = ∑ 𝑛 𝑎𝑛 𝑥 𝑛−1
′
𝑛=1
∞
𝑦 = ∑ 𝑛 (𝑛 − 1)𝑎𝑛 𝑥 𝑛−2
′′
𝑛=2
Step 3: Need to shift the first series down by 2 to get common exponents in all the series
∞ ∞ ∞
Step 4: Combine the series by stripping out the n = 0 terms from both the first and third series
∞ ∞ ∞
𝑛=1 𝑛=0
160
Step 6: Solve the first as well as the recurrence relation.
𝑎0
𝑛=0 𝑎2 = −
2
(2𝑛 − 1)𝑎𝑛
𝑛 = 1, 2, 3, … 𝑎𝑛+2 =
(𝑛 + 2)(𝑛 + 1)
Step 7: Start plugging in values of n.
𝑎0
𝑛=0 𝑎2 = −
2
𝑎1
𝑛=1 𝑎3 =
6
3𝑎2 𝑎2 𝑎0
𝑛=2 𝑎4 = = =−
12 4 8
5𝑎3 𝑎3 𝑎1
𝑛=3 𝑎5 = = =
20 4 24
Note: Can choose any arbitrary constants for a0 and a1
𝑦(𝑥) = ∑ 𝑎𝑛 𝑥 𝑛 = 𝑎0 + 𝑎1 𝑥 + 𝑎2 𝑥 2 + 𝑎3 𝑥 3 + ⋯
𝑛=0
𝑎0 2 𝑎1 3 𝑎0 𝑎1
= 𝑎0 + 𝑎1 𝑥 + (− )𝑥 + 𝑥 + (− ) 𝑥 4 + 𝑥 5 + ⋯
2 6 8 24
Step 9: Collect up the terms that contain the same coefficient, factor the coefficient out and write the
results as a new series
𝑥2 𝑥4 𝑥3 𝑥5
𝑦(𝑥) = 𝑎0 [1 − − + ⋯ ] + 𝑎1 [ 𝑥 + + + ⋯]
2 8 6 24
Step 10: Applying the initial conditions gives values for a0 and a1
𝑦(0) = 1 ⇒ 𝑎0 = 1
𝑦′(0) = 1 ⇒ 𝑎1 = 1
161
Solutions About Singular Points
The power series method for solving linear differential equations with variable coefficients no longer
works when solving the differential equation about a singular point. It appears that some features of the
solutions of such equations of the most importance for applications are largely determined by their
behavior near their singular points. Frobenius method is usually used to solve the differential equation
about a regular singular point. This method does not always yield two infinite series solutions. When only
one solution is found, a certain formula can be used to get the second solution.
For example, x = 0 is the only singular point of the Bessel equation of order n,
𝑥 2 𝑦 ′′ + 𝑥𝑦 ′ + (𝑥 2 − 𝑛2 )𝑦 = 0
whereas the Legendre equation of order n,
(1 − 𝑥 2 )𝑦 ′′ − 2𝑥𝑦 ′ + 𝑛(𝑛 + 1)𝑦 = 0
has two singular points x = -1 and x = 1.
Note: Usually, only the case in which x = 0 is a singular point of Equation (7) is considered. A differential
equation having x = a as a singular point is easily transformed by the substitution t = x – a into
one having a corresponding singular point at 0.
𝑦(𝑥) = ∑ 𝑐𝑛 𝑥 𝑛
162
So the straightforward method of power series fails in this case.
Example:
Find the singular point(s) for the differential equation
(𝑥 2 − 4)2 𝑦 ′′ + 3(𝑥 − 2)𝑦 ′ + 5𝑦 = 0
Answer
Divide the equation with
(𝑥 2 − 4)2 = (𝑥 − 2)2 (𝑥 + 2)2
and reduce the coefficients to the lowest terms, produce
3 5
𝑃(𝑥) = and 𝑄(𝑥) =
(𝑥 − 2)(𝑥 + 2)2 (𝑥 − 2)2 (𝑥 + 2)2
163
3 5
𝑝(𝑥) = (𝑥 − 2)𝑃(𝑥) = and 𝑞(𝑥) = (𝑥 − 2)2 𝑄(𝑥) =
(𝑥 + 2)2 (𝑥 + 2)2
are analytic at x -2.
(ii) Now since the factor x - (-2) = x + 2 appears to the second power in the denominator of P(x), we
can conclude immediately that x = -2 is an irregular singular point of the equation. This also
follows from the fact that
3
𝑝(𝑥) = (𝑥 + 2)𝑃(𝑥) =
(𝑥 − 2)(𝑥 + 2)
is not analytic at x = -2.
If x = x0 is a singular point of the differential equation (8), then there exists at least one solution of the
form
∞ ∞
𝑦(𝑥) = (𝑥 − 𝑥0 )𝑟 ∑ 𝑐𝑛 (𝑥 − 𝑥0 )𝑛 = ∑ 𝑐𝑛 (𝑥 − 𝑥0 )𝑛+𝑟
𝑛=0 𝑛=0
where the number r is a constant to be determined. The series will converge at least on some interval
0 < x – x0 < R.
where the exponent r may be any real or complex number and r is chosen so that a0 0.
The equation also has a second solution such that these two solutions are linearly independent that may
be similar to (10) with different r and different coefficients or may contain a logarithmic term.
164
Solution Procedure:
(i) Write (8) in the following more convenient form
𝑥 2 𝑦 ′′ + 𝑥𝑏(𝑥)𝑦 ′ + 𝑐(𝑥)𝑦 = 0 (11)
𝑏(𝑥) = ∑ 𝑏𝑛 𝑥𝑛 = 𝑏0 + 𝑏1 𝑥 + 𝑏2 𝑥 2 + ⋯
𝑛=0
∞
𝑐(𝑥) = ∑ 𝑐𝑛 𝑥𝑛 = 𝑐0 + 𝑐1 𝑥 + 𝑐2 𝑥 2 + ⋯
𝑛=0
When n = 0
𝑥 𝑟 [𝑟(𝑟 − 1)𝑎0 + 𝑟𝑎0 𝑏0 + 𝑎0 𝑐0 ] = 0
𝑥𝑟 ≠ 0 ⇒ [𝑟(𝑟 − 1)𝑎0 + 𝑟𝑏0 + 𝑐0 ]𝑎0 = 0
𝑎0 ≠ 0 ⇒ 𝑟(𝑟 − 1)𝑎0 + 𝑟𝑏0 + 𝑐0 = 0
𝑟 2 + (𝑏0 − 1)𝑟 + 𝑐0 = 0 (12)
165
Equation (12) is called an indicial equation of the differential equation (11), and its two roots
(possibly equal) are exponents of the differential equation at the regular singular point x = 0.
Equation (12) also can be written as
𝑟(𝑟 − 1) + 𝑏0 𝑟 + 𝑐0 = 0 (13)
(v) Find the indicial roots r1, r2 for the indicial equation
If (9) is to be a solution of the differential equation in (10), the exponent r must be one of the roots
r1 and r2 of the indicial equation in (13). If r1 r2, it follows that there are two possible Frobenius
series solutions, whereas if r1 = r2 there is only one possible Frobenius series solution; the second
solution cannot be a Frobenius series.
The exponents r1 and r2 in the possible Frobenius series solutions are determined (using the
indicial equation) by the values b0 = b(0) and c0 = c(0). In practice, particularly when the
coefficients in the differential equation in the original form in (8) are polynomials, the simplest
way of finding b0 and c0 is often to write the equation in the form
𝑏0 + 𝑏1 𝑥 + 𝑏2 𝑥 2 + ⋯ ′ 𝑐0 + 𝑐1 𝑥 + 𝑐2 𝑥 2 + ⋯
𝑦 ′′ + 𝑦 + 𝑦=0
𝑥 𝑥2
Then inspection of the series that appear in the two numerators reveals the constants b0
and c0.
Example
Find the exponents in the possible Frobenius series solutions of the equation
2𝑥 2 (1 + 𝑥)𝑦 ′′ + 3𝑥(1 + 𝑥)3 𝑦 ′ − (1 − 𝑥 2 )𝑦 = 0
Answer:
Divide each term by 2𝑥 2 (1 + 𝑥) to recast the differential equation in the form
3 1
′′ 2
(1 + 2𝑥 + 𝑥 2 ) ′ −
2
(1 − 𝑥)
𝑦 + 𝑦 + 𝑦=0
𝑥 𝑥2
3 1
and get 𝑏0 = 2 and 𝑐0 = − 2 . Hence the indicial equation is
3 1 1 1 1
𝑟(𝑟 − 1) + 𝑟 − = 𝑟 2 + 𝑟 − = (𝑟 + 1) (𝑟 − ) = 0
2 2 2 2 2
1
with roots 𝑟1 = 2 and 𝑟2 = −1. The two possible Frobenius series solutions are then in the forms
∞ ∞
1⁄2 𝑛 −1
𝑦1 (𝑥) = 𝑥 ∑ 𝑎𝑛 𝑥 and 𝑦2 (𝑥) = 𝑥 ∑ 𝑎𝑛 𝑥 𝑛
𝑛=0 𝑛=0
166
Review Materials
1. Taylor Series
If f(x) is an infinitely differentiable function then the Taylor Series of f(x) about x = x0 is,
∞
𝑓 (𝑛) (𝑥0 )
𝑓(𝑥) = ∑ (𝑥 − 𝑥0 )𝑛
𝑛!
𝑛=0
Recall that
𝑓 (0) (𝑥) = 𝑓(𝑥)
𝑓 (𝑛) (𝑥) = 𝑛𝑡ℎ derivative of 𝑓(𝑥)
Example
Determine the Taylor series for 𝑓(𝑥) = 𝑒 𝑥 about x = 0
Solution
This is probably one of the easiest functions to find the Taylor series for. Recall that,
𝑓 (𝑛) (𝑥) = 𝑒 𝑥 𝑛 = 0, 1, 2, ⋯
and get
𝑓 (𝑛) (0) = 1 𝑛 = 0, 1, 2, ⋯
The Taylor series for this example is then
∞
𝑥
𝑥𝑛
𝑒 =∑
𝑛!
𝑛=0
Example
Determine the Taylor series for 𝑓(𝑥) = 𝑒 𝑥 about x = -4
Solution
This problem is virtually identical to the previous problem. In this case we just need to notice that,
𝑓 (𝑛) (−4) = 𝑒 −4 𝑛 = 0, 1, 2, ⋯
The Taylor series for this example is then,
∞
𝑥
𝑒 −4
𝑒 =∑ (𝑥 + 4)𝑛
𝑛!
𝑛=0
167
Definition
A function, f(x), is called analytic at x = a if the Taylor series for f(x) about x = a has a positive radius
of convergence and converges to f(x).
168
ENGINEERING APPLICAIONS OF
DIFFERENTIAL EQUATION
WEEK 14: ENGINEERING APPLICATIONS OF DIFFERENTIAL EQUATION
14.1 A LIQUID SYSTEM
Figure 1 shows a tank of liquid. The tank has a constant cross-sectional area A. The liquid can flow out of
the tank through a valve near the base. As it does, the height or head, h, of liquid in the tank will reduce.
Let q be the rate at which liquid flows out of the tank. Under certain conditions the rate outflow is
proportional to the head, so that q = kh where k is a constant of proportionality. Situation like this arise
frequently in chemical engineering industry. Obtain a differential equation for this system or equivalently,
come out with the mathematical model for the physical system.
Area, A
q = kh
q
Figure 1. Modelling a liquid system
The expression for the volume V of liquid in the tank at any time.
𝑉 =𝐴×ℎ
The volume of liquid in the tank changes because liquid is flowing out. Hence,
the rate at which this volume changes = rate of flow in – rate of flow out
This is the law of conservation of mass. The rate of change of volume is
𝑑𝑉
𝑑𝑡
There is no flow into the tank and liquid flows out at a rate q. Hence
𝑑𝑉
= −𝑞
𝑑𝑡
𝑑ℎ
But, 𝑉 = 𝐴ℎ and A is constant, so the rate of change of volume is 𝐴 𝑑𝑡
169
Therefore
𝑑𝑉 𝑑ℎ
=𝐴 = −𝑞
𝑑𝑡 𝑑𝑡
Also, 𝑞 = 𝑘ℎ, so
𝑑ℎ
𝐴 = −𝑘ℎ
𝑑𝑡
or
𝐴ℎ′ + 𝑘ℎ = 0
This is a first order differential equation with dependent variable h and independent variable t. It is linear
and has constant coefficients. The unknown function to seek (the solution) is h(t). Solve the equation to
find the head, h, at any time, t.
A tank contains 40 gallons of a solution composed of 90 percent water and 10 percent alcohol. A second
solution containing 50 percent water and 50 percent alcohol is added to the tank at a rate of 4
gallon/minute. As the second solution being added, the tank is being drained at the rate of 4
gallon/minute, as shown in Figure 2. Assuming the solution in the tank is stirred constantly, how much
alcohol is in the tank after t minutes, how much alcohol is in the tank after 10 minutes?
170
where 2 represents the 2 gallons of alcohol entering each minute in the 50% solution.
In standard form,
1
𝑦′ + 𝑦= 2
10
is a first order linear differential equation. The unknown function to seek (the solution) is y(t). Solve the
equation to find the number of gallons of alcohol, y, at any time, t.
Solution
1
To solve the linear equation, let 𝑝(𝑥) =
10
Obtain
1 𝑡
∫ 𝑝(𝑡)𝑑𝑡 = ∫ 𝑑𝑡 =
10 10
Thus, the integrating factor is
𝑡⁄
𝑢(𝑡) = 𝑒 ∫ 𝑝(𝑡)𝑑𝑡 = 𝑒 10
𝑡 𝑡⁄
= 𝑒 − ⁄10 (20𝑒 10 + 𝐶)
𝑡
= 20 + 𝐶𝑒 − ⁄10
Initial conditions, since y = 4 when t = 0,
4 = 20 + 𝐶 ⇒ 𝐶 = −16
The particular solution is
𝑡
𝑦 = 20 − 16𝑒 − ⁄10
Finally, when t = 10, the amount of alcohol in the tank is
10⁄
𝑦 = 20 − 16𝑒 − 10 = 20 − 16𝑒 −1 ≈ 14.11 gallons
171
14.3 AN LCR CIRCUIT
Figure 3 shows an LCR circuit. This is a circuit comprising an inductor of inductance L, a capacitor of
capacitance C, and a resistor of resistance R placed in a series. When a constant voltage source, V, is
applied, it can be shown that the current, i, through the circuit satisfies the differential equation
𝑑2 𝑖 𝑑𝑖 1
𝐿 2+𝑅 + 𝑖 =0
𝑑𝑡 𝑑𝑡 𝐶
or
1
𝐿 𝑖 ′′ + 𝑅 𝑖 ′ + 𝑖=0
𝐶
This equation can be derived using Kirchhoff’s voltage law, the individual laws for each component.
Because L, R, and C are constants, this is a constant coefficient equation. It is linear and second order.
The unknown function to seek (the solution) is i(t). Solve the equation to find the current in the circuit, i,
at any time, t.
(b)
(a)
172
Consider the motion of an object with mass m at the end of the spring that is either vertical or horizontal
on a level surface. Using Hooke’s Law, which says that if the spring is stretched (or compressed) x units
from its natural length, then it exerts a force that is proportional to x
restoring force = −𝑘𝑥
where k is a positive constant (called the spring constant). If any external resisting forces (due to air
resistance or friction) are ignored, then by Newton’s Second Law, F = ma,
𝑑2 𝑥
𝑚 = −𝑘𝑥
𝑑𝑡 2
or
𝑑2 𝑥
𝑚 + 𝑘𝑥 = 0 (1)
𝑑𝑡 2
This is a second order linear differential equation.
Solution:
The auxiliary equation is
𝑚𝑟 2 + 𝑘 = 0
with roots
𝑘
𝑟 = ±𝜔𝑖 where 𝜔=√
𝑚
𝑘
𝜔=√ (frequency)
𝑚
𝑐1 𝑐2
cos 𝛿 = sin 𝛿 = − (phase angle)
𝐴 𝐴
Exercise 1
173
A spring with mass of 2kg has natural length 0.5 m. A force of 25.6 N is required to maintain it stretched
to a length of 0.7 m. If the spring is stretched to a length of 0.7 m and then released with initial velocity 0,
find the position of the mass at any time t.
Solution
From Hooke’s Law, the force required to stretch the spring is
𝑘(0.2) = 25.6
25.6
𝑘= = 128
0.2
Using this value of the spring constant k, together with m = 2 in Equation 1
𝑑2 𝑥
2 + 128𝑥 = 0
𝑑𝑡 2
The solution of this equation is
𝑥(𝑡) = 𝑐1 cos 8𝑡 + 𝑐2 sin 8𝑡 (2)
The initial condition is given as 𝑥(0) = 0.2. Hence from Equation (2), 𝑥(0) = 𝑐1 = 0.2.
Differentiating Equation (2)
𝑥′(𝑡) = −8𝑐1 cos 8𝑡 + 8𝑐2 sin 8𝑡
Since the initial velocity is given as 𝑥′(0) = 0, 𝑐2 = 0
So, the solution is
𝑥(𝑡) = 0.2 cos 8𝑡
Consider the motion of a spring that is subject to a frictional force (in the case of horizontal spring) or a
damping force (in the case where a vertical spring moves through a fluid) An example is the damping
force supplied by a shock absorber in a car or a bicycle.
174
Assume that the damping force is proportional to the velocity of the mass and acts in the direction
opposite to the motion. Thus,
𝑑𝑥
damping force = −𝑐
𝑑𝑡
where c is a positive constant, called the damping constant. Thus, Newton’s Second Law gives
𝑑2 𝑥 𝑑𝑥
𝑚 2
= restoring force + damping force = −𝑘𝑥 − 𝑐
𝑑𝑡 𝑑𝑡
or
𝑑2 𝑥 𝑑𝑥
𝑚 2
+𝑐 + 𝑘𝑥 = 0 (3)
𝑑𝑡 𝑑𝑡
This is a second order linear differential equation.
Solution:
The auxiliary equation is
𝑚𝑟 2 + 𝑐𝑟 + 𝑘 = 0
with roots
−𝑐 + √𝑐 2 − 4𝑚𝑘 −𝑐 − √𝑐 2 − 4𝑚𝑘
𝑟1 = 𝑟2 = (4)
2𝑚 2𝑚
Since c, m, and k are all positive, then √𝑐 2 − 4𝑚𝑘 < 0, and the roots r1 and r2 given by Equations (3) must
be both negative. This shows that 𝑥 → 0 as 𝑡 → ∞. Typical graphs of x as a function of t are shown in
Figure 6. Notice oscillations do not occur. This is because 𝑐 2 > 4𝑚𝑘 means that there is a strong
damping force (high viscosity oil or grease) compared with a weak spring or small mass.
175
Figure 6. Typical graphs for overdamping case
and a typical graph is shown in Figure 7. It is similar to Case I, but the damping is just sufficient to suppress
vibrations. Any decrease in the viscosity of fluid leads to the vibrations of the following case.
√4𝑚𝑘 − 𝑐 2
𝜔=
2𝑚
The solution is given by
176
𝑐
−( )𝑡
𝑥=𝑒 2𝑚 (𝑐1 cos 𝜔𝑡 + 𝑐2 sin 𝜔𝑡)
𝑐 𝑐
−( )𝑡 −( )𝑡
There are oscillations that are damped by the factor 𝑒 2𝑚 . Since c > 0 and m > 0, then 𝑒 2𝑚 < 0 so
𝑐
−( )𝑡
𝑒 2𝑚 → 0 as 𝑡 → ∞ . This implies that 𝑥 → 0 as 𝑡 → ∞ ; that is, the motion decays to 0 as time
increases. A typical graph is shown in Figure 8
Exercise 2
Suppose that the string in Exercise 1 is immersed in a fluid with damping constant c = 40. Find the position
of the mass at any time t if it starts from the equilibrium position and is given a push to start it with an
initial velocity of 0.6 m/s.
Solution
The mass is m = 2 and the spring constant is k = 128, so the differential equation becomes
𝑑2 𝑥 𝑑𝑥
2 2
+ 40 + 128𝑥 = 0
𝑑𝑡 𝑑𝑡
or
𝑑2 𝑥 𝑑𝑥
+ 20 + 64𝑥 = 0
𝑑𝑡 2 𝑑𝑡
The auxiliary equation is
𝑟 2 + 20𝑟 + 64 = (𝑟 + 4)(𝑟 + 16) = 0
with roots -4 and -16, so the motion is overdamped and the solution is
𝑥(𝑡) = 𝑐1 𝑒 −4𝑡 + 𝑐2 𝑒 −16𝑡
Given 𝑥(0) = 0, so 𝑐1 + 𝑐2 = 0. Differentiating,
𝑥 ′ (𝑡) = −4𝑐1 𝑒 −4𝑡 − 16𝑐2 𝑒 −16𝑡
So
177
𝑥 ′ (0) = −4𝑐1 − 16𝑐2 = 0.6
Since 𝑐1 = −𝑐2, this gives 12𝑐1 = 0.6 or 𝑐1 = 0.05. Therefore
𝑥(𝑡) = 0.05(𝑒 −4𝑡 − 𝑒 −16𝑡 )
Suppose that, in addition to the restoring force and the damping force, the motion of the spring is
affected by an external force F(t). The Newton’s Second Law gives
𝑑2 𝑥 𝑑𝑥
𝑚 2
= restoring force + damping force + external force = −𝑘𝑥 − 𝑐 + 𝐹(𝑡)
𝑑𝑡 𝑑𝑡
Thus, instead of the homogeneous Equation (3), the motion of the spring is now governed by the
following nonhomogeneous differential equation.
𝑑2 𝑥 𝑑𝑥
𝑚 2
+𝑐 + 𝑘𝑥 = 𝐹(𝑡) (5)
𝑑𝑡 𝑑𝑡
The motion of the spring can be determined by the methods such as the method of undetermined
coefficients, the method of variation parameters, etc.
A commonly occurring type of external force is a periodic force function
𝐹(𝑡) = 𝐹0 cos 𝜔0 𝑡
𝑘
where 𝜔0 ≠ 𝜔 = √𝑚
In this case, and in the absence of a damping force (c = 0), using method of undetermined coefficients
show that
𝐹0
𝑥 = 𝑐1 cos 𝜔𝑡 + 𝑐2 sin 𝜔𝑡 + cos 𝜔0 𝑡 (6)
𝑚(𝜔 2 − 𝜔02 )
If 𝜔0 = 𝜔, then the applied frequency reinforces the natural frequency and the result is vibrations of large
amplitude. This is the phenomenon of resonance.
178
14.7 ELECTRIC CIRCUIT
Analyze the electric circuit shown in Figure 9. It contains an electromotive force E (supplied by a battery
or generator), a resistor R, an inductor L, and a capacitor C, in series. If the charge on the capacitor at
time t is Q = Q(t), then the current, I, is the rate of change of Q with respect to t,
𝑑𝑄
𝐼=
𝑑𝑡
It is known from physics that the voltage drops across the resistor, inductor, and capacitor are
𝑑𝐼 𝑄
𝑅𝐼 𝐿
𝑑𝑡 𝐶
respectively. Kirchhoff’s voltage law says that the sum of these voltage drops is equal to the supplied
voltage
𝑑𝐼 𝑄
𝐿 + 𝑅𝐼 + = 𝐸(𝑡) (7)
𝑑𝑡 𝐶
𝑑𝑄
Since 𝐼 = , Equation 7 becomes
𝑑𝑡
𝑑2 𝑄 𝑑𝑄 1
𝐿 2
+𝑅 + 𝑄 = 𝐸(𝑡) (8)
𝑑𝑡 𝑑𝑡 𝐶
which is a second order linear differential equation with constants coefficients. If the charge Q0 and the
current I0 are known at time 0, then the initial conditions are
𝑄(0) = 𝑄0 𝑄 ′ = 𝐼(0) = 𝐼0
The initial value problem can be solved by such as the method of undetermined coefficients, the
method of variation parameters, etc.
The differential equation for the current can be obtained by differentiating Equation (*) with respect to
𝑑𝑄
t and using 𝐼 = 𝑑𝑡
179
𝑑2 𝐼 𝑑𝐼 1
𝐿 2
+ 𝑅 + 𝐼 = 𝐸′(𝑡)
𝑑𝑡 𝑑𝑡 𝐶
Exercise 3
Find the charge and current at time t in the circuit of Figure 9 if R = 40 , L = 1 H, C = 16 × 10-4 F, E(t) =
100 cos 10t, and the initial charge and current are both 0.
Solution
𝑑2 𝑄 𝑑𝑄
2
+ 40 + 625𝑄 = 100 cos 10𝑡 (9)
𝑑𝑡 𝑑𝑡
−40 ± √−900
𝑟= = −20 ± 15𝑖
2
The solution of the complementary equation is
Then
(525𝐴 + 400𝐵) cos 10𝑡 + (−400𝐴 + 525𝐵) sin 10𝑡 = 100 cos 10𝑡
Equating coefficients,
or
The solution is
180
84 64
𝐴= 𝐵=
697 697
and the particular solution is
1
𝑄𝑝 (𝑡) = (84 cos 10𝑡 + 64 sin 10𝑡)
697
The general solution is
4
𝑄(𝑡) = 𝑄𝑐 (𝑡) + 𝑄𝑝 (𝑡) = 𝑒 −20𝑡 (𝑐1 cos 15𝑡 + 𝑐2 sin 15𝑡) + (21 cos 10𝑡 + 16 sin 10𝑡)
697
Imposing the initial condition Q(0) = 0
84 84
𝑄(0) = 𝑐1 + =0 or 𝑐1 = −
697 697
To impose the other initial condition, first differentiate to find the current
𝑑𝑄 40
𝐼= = 𝑒 −20𝑡 [(−20𝑐1 + 15𝑐2 ) cos 15𝑡 + (−15𝑐1 − 20𝑐2 ) sin 15𝑡] + (−21 sin 10𝑡 + 16 cos 10𝑡)
𝑑𝑡 697
640 464
𝐼(0) = −20𝑐1 + 15𝑐2 + =0 or 𝑐2 = −
697 2091
Thus, the formula for charge is
4 𝑒 −20𝑡
𝑄(𝑡) = [ (−63 cos 15𝑡 − 116 sin 15𝑡) + (21 cos 10𝑡 + 16 sin 10𝑡)]
697 3
1
𝐼(𝑡) = [𝑒 −20𝑡 (−1920 cos 15𝑡 + 13.060 sin 15𝑡) + 120(−21 sin 10𝑡 + 16 cos 10𝑡)]
2091
Note 1:
The solution for Q(t) consists of two parts. Since 𝑒 −20𝑡 → 0 as 𝑡 → ∞ and both cos 15t and sin 15t are
bounded functions
4
𝑄𝑐 (𝑡) = 𝑒 −20𝑡 (−63 cos 15𝑡 − 116 sin 15𝑡) → 0 as 𝑡 → ∞
2091
So, for large values of t,
4
𝑄(𝑡) ≈ 𝑄𝑝 (𝑡) = (21 cos 10𝑡 + 16 sin 10𝑡)
697
181
and, for this reason, Qp(t) is called the steady-state solution. Figure 10 shows how the graph of the steady
state solution compares with the graph of Q.
Note 2:
Physical situations versus mathematically defined for the spring system and electric circuit are given in
Table 1.
182
APPENDIX 11.1 MATHEMATICAL MODELING AND ENGINEERING PROBLEM SOLVING
Mathematical model plays an important role in engineering problem solving. The engineering problem
solving process in parachute design is illustrated as follows:
To develop the mathematical model for calculating the velocity of a parachute, we can use the
knowledge from the existing theory/physical law (e.g. Newton’s Law) or understand the problem by
empirical means (e.g. by observation and experiment). In this case, Newton’s 2nd Law is applied, where
∑ 𝐹(𝑡) = 𝑚𝑎(𝑡). A free-body diagram is drawn as below. Assume that the drag force due to air resistance
is proportional to the falling velocity, the downward motion is positive, and 𝑔 = 9.81𝑚𝑠 −2. Then, we are
able to develop the mathematical model.
183
The mathematical modelling of the falling parachutist’s velocity:
𝑑𝑣(𝑡)
∑ 𝐹(𝑡) = 𝑚𝑎(𝑡) = 𝑚 .
𝑑𝑡
𝑑𝑣(𝑡)
𝐹𝑔𝑟𝑎𝑣𝑖𝑡𝑦 + 𝐹𝑑𝑟𝑎𝑔 = 𝑚
𝑑𝑡
𝑑𝑣(𝑡) 𝐹𝑔𝑟𝑎𝑣𝑖𝑡𝑦 + 𝐹𝑑𝑟𝑎𝑔 𝑐𝑣(𝑡)
𝑅𝑎𝑡𝑒 𝑜𝑓 𝑣𝑒𝑙𝑜𝑐𝑖𝑡𝑦 𝑐ℎ𝑎𝑛𝑔𝑒 = = =𝑔−
𝑑𝑡 𝑚 𝑚
Rearrange it, we found that the mathematical model is in the form of first order linear
nonhomogeneous Ordinary Differential Equation (ODE):
The classification of the order, linear vs nonlinear, homogeneous vs nonhomogeneous will be covered
later.
184
Table 12.1: Elements of mathematical model and its example of falling parachutist problem.
With this solution, it helps us to gain intuition about what to expect on the behaviour of the
examined system. For example, what are the suitable mass and damping to reduce the impact velocity
of the parachutist to the ground? This approach is known as the result interpretation and analysis. With
an accurate mathematical model, we can predict the performance of the design without testing with the
real subject. Avoiding the unnecessary physical cycles of ‘modify-and-test’ would save time and money.
In addition, numerical method is used to solve a complicated mathematical model that can’t be
solved by using analytical method. The desired solution obtained by numerical method- Euler method is
𝑐
given as follows: 𝑣(𝑡𝑖+1 ) = 𝑣(𝑡𝑖 ) + [𝑔 − 𝑚 𝑣(𝑡𝑖 )]( 𝑡𝑖+1 − 𝑡𝑖 ). This is for your extra information and you
will learn the numerical method in advanced mathematic class.
185
APPENDIX 11.2 CONSERVATION LAWS AND ENGINEERING
When we employed the Newton’s law to develop a force balance equation for the falling parachutist, i.e.
𝑑𝑣(𝑡) 𝑚𝑔−𝑐𝑣(𝑡)
𝑑𝑡
= 𝑚
, we can eventually boil down it to a simple equation:
𝑑𝑣(𝑡)
(i) Transient problem: [Change=increase-decrease] or [ 𝑑𝑡
≠ 0]
Although simple, it embodies one of the most fundamental ways in which conservation laws are
used in engineering – that is, to predict changes with respect to time. In this case, it is recognized as time-
variant (or transient) problem. (For example, the falling velocity of the parachutist with 𝑚 = 68.1𝑘𝑔, 𝑐 =
12.5𝑘𝑔/𝑠 changes with respect to time is given in Figure 11.2 (a).)
Asides from predicting changes, another way in which conservation laws are applied is for cases
where change is nonexistent. In this case, it is recognized as time-invariant (or steady state) problem (For
example, we would like to know when will the falling velocity becomes constant, i.e. the terminal
velocity).
𝑑𝑣(𝑡)
(ii) Steady state problem: [Change=increase-decrease=0] or [ 𝑑𝑡
= 0]
𝑑𝑣(𝑡)
For falling parachutist case, steady-state conditions would correspond to the case where 𝑑𝑡
=
𝑚𝑔−𝑐𝑣(𝑡) 𝑔𝑚
𝑚
= 0. Figure 11.2 is plotted by using the solution, 𝑣(𝑡) = 𝑐
(1 − 𝑒 −(𝑐⁄𝑚)𝑡 ). From the Figure 11.2
below, it shows that the falling parachutist’s velocity keeps increasing and thus varies with time initially
(transient case). When the velocity reaches to a point where the 𝐹𝑔𝑟𝑎𝑣𝑖𝑡𝑦 , (𝑚𝑔) is equal to 𝐹𝑑𝑟𝑎𝑔 , (𝑐𝑣),
𝑑𝑣(𝑡)
there is no more increase in velocity afterward due to 𝑑𝑡
= 0.
186
Figure 11.2: (a) Transient and (b) steady state cases for falling parachutist problem.
𝑔𝑚
From Figure 11.2(b), the terminal velocity and terminal time can be computed as 𝑣𝑡𝑒𝑟𝑚𝑖𝑛𝑎𝑙 = 𝑐
𝑚 𝑐
and 𝑡𝑡𝑒𝑟𝑚𝑖𝑛𝑎𝑙 = − 𝑐
ln |1 − 𝑔𝑚 (𝑣𝑡𝑒𝑟𝑚𝑖𝑛𝑎𝑙 )|. Note that velocity will increase when the first derivative,
𝑑𝑣(𝑡) 𝑑𝑣(𝑡) 𝑑𝑣(𝑡)
> 0 , decrease when < 0, and remain unchanged when = 0.
𝑑𝑡 𝑑𝑡 𝑑𝑡
187